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# %%
import argparse, os
import torch
import requests
import torch.nn as nn
import torch.nn.functional as F
from PIL import Image
from io import BytesIO
from tqdm.auto import tqdm
from matplotlib import pyplot as plt
from torchvision import transforms as tfms
from diffusers import (
StableDiffusionPipeline,
DDIMScheduler,
DiffusionPipeline,
StableDiffusionXLPipeline,
)
from diffusers.image_processor import VaeImageProcessor
import torch
import torch.nn as nn
import torchvision
import torchvision.transforms as transforms
from torchvision.utils import save_image
import argparse
import PIL.Image as Image
from torchvision.utils import make_grid
import numpy
from diffusers.schedulers import DDIMScheduler
import torch.nn.functional as F
from models import attn_injection
from omegaconf import OmegaConf
from typing import List, Tuple
import omegaconf
import utils.exp_utils
import json
device = "cuda"
def _get_text_embeddings(prompt: str, tokenizer, text_encoder, device):
# Tokenize text and get embeddings
text_inputs = tokenizer(
prompt,
padding="max_length",
max_length=tokenizer.model_max_length,
truncation=True,
return_tensors="pt",
)
text_input_ids = text_inputs.input_ids
with torch.no_grad():
prompt_embeds = text_encoder(
text_input_ids.to(device),
output_hidden_states=True,
)
pooled_prompt_embeds = prompt_embeds[0]
prompt_embeds = prompt_embeds.hidden_states[-2]
if prompt == "":
negative_prompt_embeds = torch.zeros_like(prompt_embeds)
negative_pooled_prompt_embeds = torch.zeros_like(pooled_prompt_embeds)
return negative_prompt_embeds, negative_pooled_prompt_embeds
return prompt_embeds, pooled_prompt_embeds
def _encode_text_sdxl(model: StableDiffusionXLPipeline, prompt: str):
device = model._execution_device
(
prompt_embeds,
pooled_prompt_embeds,
) = _get_text_embeddings(prompt, model.tokenizer, model.text_encoder, device)
(
prompt_embeds_2,
pooled_prompt_embeds_2,
) = _get_text_embeddings(prompt, model.tokenizer_2, model.text_encoder_2, device)
prompt_embeds = torch.cat((prompt_embeds, prompt_embeds_2), dim=-1)
text_encoder_projection_dim = model.text_encoder_2.config.projection_dim
add_time_ids = model._get_add_time_ids(
(1024, 1024), (0, 0), (1024, 1024), torch.float16, text_encoder_projection_dim
).to(device)
# repeat the time ids for each prompt
add_time_ids = add_time_ids.repeat(len(prompt), 1)
added_cond_kwargs = {
"text_embeds": pooled_prompt_embeds_2,
"time_ids": add_time_ids,
}
return added_cond_kwargs, prompt_embeds
def _encode_text_sdxl_with_negative(
model: StableDiffusionXLPipeline, prompt: List[str]
):
B = len(prompt)
added_cond_kwargs, prompt_embeds = _encode_text_sdxl(model, prompt)
added_cond_kwargs_uncond, prompt_embeds_uncond = _encode_text_sdxl(
model, ["" for _ in range(B)]
)
prompt_embeds = torch.cat(
(
prompt_embeds_uncond,
prompt_embeds,
)
)
added_cond_kwargs = {
"text_embeds": torch.cat(
(added_cond_kwargs_uncond["text_embeds"], added_cond_kwargs["text_embeds"])
),
"time_ids": torch.cat(
(added_cond_kwargs_uncond["time_ids"], added_cond_kwargs["time_ids"])
),
}
return added_cond_kwargs, prompt_embeds
# Sample function (regular DDIM)
@torch.no_grad()
def sample(
pipe,
prompt,
start_step=0,
start_latents=None,
intermediate_latents=None,
guidance_scale=3.5,
num_inference_steps=30,
num_images_per_prompt=1,
do_classifier_free_guidance=True,
negative_prompt="",
device=device,
):
negative_prompt = [""] * len(prompt)
# Encode prompt
if isinstance(pipe, StableDiffusionPipeline):
text_embeddings = pipe._encode_prompt(
prompt,
device,
num_images_per_prompt,
do_classifier_free_guidance,
negative_prompt,
)
added_cond_kwargs = None
elif isinstance(pipe, StableDiffusionXLPipeline):
added_cond_kwargs, text_embeddings = _encode_text_sdxl_with_negative(
pipe, prompt
)
# Set num inference steps
pipe.scheduler.set_timesteps(num_inference_steps, device=device)
# Create a random starting point if we don't have one already
if start_latents is None:
start_latents = torch.randn(1, 4, 64, 64, device=device)
start_latents *= pipe.scheduler.init_noise_sigma
latents = start_latents.clone()
latents = latents.repeat(len(prompt), 1, 1, 1)
# assume that the first latent is used for reconstruction
for i in tqdm(range(start_step, num_inference_steps)):
latents[0] = intermediate_latents[(-i + 1)]
t = pipe.scheduler.timesteps[i]
# Expand the latents if we are doing classifier free guidance
latent_model_input = (
torch.cat([latents] * 2) if do_classifier_free_guidance else latents
)
latent_model_input = pipe.scheduler.scale_model_input(latent_model_input, t)
# Predict the noise residual
noise_pred = pipe.unet(
latent_model_input,
t,
encoder_hidden_states=text_embeddings,
added_cond_kwargs=added_cond_kwargs,
).sample
# Perform guidance
if do_classifier_free_guidance:
noise_pred_uncond, noise_pred_text = noise_pred.chunk(2)
noise_pred = noise_pred_uncond + guidance_scale * (
noise_pred_text - noise_pred_uncond
)
latents = pipe.scheduler.step(noise_pred, t, latents).prev_sample
# Post-processing
images = pipe.decode_latents(latents)
images = pipe.numpy_to_pil(images)
return images
# Sample function (regular DDIM), but disentangle the content and style
@torch.no_grad()
def sample_disentangled(
pipe,
prompt,
start_step=0,
start_latents=None,
intermediate_latents=None,
guidance_scale=3.5,
num_inference_steps=30,
num_images_per_prompt=1,
do_classifier_free_guidance=True,
use_content_anchor=True,
negative_prompt="",
device=device,
):
negative_prompt = [""] * len(prompt)
vae_decoder = VaeImageProcessor(vae_scale_factor=pipe.vae.config.scaling_factor)
# Encode prompt
if isinstance(pipe, StableDiffusionPipeline):
text_embeddings = pipe._encode_prompt(
prompt,
device,
num_images_per_prompt,
do_classifier_free_guidance,
negative_prompt,
)
added_cond_kwargs = None
elif isinstance(pipe, StableDiffusionXLPipeline):
added_cond_kwargs, text_embeddings = _encode_text_sdxl_with_negative(
pipe, prompt
)
# Set num inference steps
pipe.scheduler.set_timesteps(num_inference_steps, device=device)
# save
latent_shape = (
(1, 4, 64, 64) if isinstance(pipe, StableDiffusionPipeline) else (1, 4, 64, 64)
)
generative_latent = torch.randn(latent_shape, device=device)
generative_latent *= pipe.scheduler.init_noise_sigma
latents = start_latents.clone()
latents = latents.repeat(len(prompt), 1, 1, 1)
# randomly initialize the 1st latent for generation
latents[1] = generative_latent
num_intermediate_latents = len(intermediate_latents) if intermediate_latents is not None else 0
for i in range(start_step, num_inference_steps):
if use_content_anchor and intermediate_latents is not None:
# Ensure that the index is within bounds
if -i >= -num_intermediate_latents:
latents[0] = intermediate_latents[-i]
else:
# Handle case when the index is out of bounds
# You could use a default latent or skip this step
latents[0] = intermediate_latents[0] # Example: use the first latent
t = pipe.scheduler.timesteps[i]
# Expand the latents if we are doing classifier free guidance
latent_model_input = (
torch.cat([latents] * 2) if do_classifier_free_guidance else latents
)
latent_model_input = pipe.scheduler.scale_model_input(latent_model_input, t)
# Predict the noise residual
noise_pred = pipe.unet(
latent_model_input,
t,
encoder_hidden_states=text_embeddings,
added_cond_kwargs=added_cond_kwargs,
).sample
# Perform guidance
if do_classifier_free_guidance:
noise_pred_uncond, noise_pred_text = noise_pred.chunk(2)
noise_pred = noise_pred_uncond + guidance_scale * (
noise_pred_text - noise_pred_uncond
)
latents = pipe.scheduler.step(noise_pred, t, latents).prev_sample
# Post-processing
# images = vae_decoder.postprocess(latents)
pipe.vae.to(dtype=torch.float32)
latents = latents.to(next(iter(pipe.vae.post_quant_conv.parameters())).dtype)
latents = 1 / pipe.vae.config.scaling_factor * latents
images = pipe.vae.decode(latents, return_dict=False)[0]
images = (images / 2 + 0.5).clamp(0, 1)
# we always cast to float32 as this does not cause significant overhead and is compatible with bfloat16
images = images.cpu().permute(0, 2, 3, 1).float().numpy()
images = pipe.numpy_to_pil(images)
if isinstance(pipe, StableDiffusionXLPipeline):
pipe.vae.to(dtype=torch.float16)
return images
## Inversion
@torch.no_grad()
def invert(
pipe,
start_latents,
prompt,
guidance_scale=3.5,
num_inference_steps=50,
num_images_per_prompt=1,
do_classifier_free_guidance=True,
negative_prompt="",
device=device,
):
# Encode prompt
if isinstance(pipe, StableDiffusionPipeline):
text_embeddings = pipe._encode_prompt(
prompt,
device,
num_images_per_prompt,
do_classifier_free_guidance,
negative_prompt,
)
added_cond_kwargs = None
latents = start_latents.clone().detach()
elif isinstance(pipe, StableDiffusionXLPipeline):
added_cond_kwargs, text_embeddings = _encode_text_sdxl_with_negative(
pipe, [prompt]
) # Latents are now the specified start latents
latents = start_latents.clone().detach().half()
# We'll keep a list of the inverted latents as the process goes on
intermediate_latents = []
# Set num inference steps
pipe.scheduler.set_timesteps(num_inference_steps, device=device)
# Reversed timesteps <<<<<<<<<<<<<<<<<<<<
timesteps = list(reversed(pipe.scheduler.timesteps))
for i in range(num_inference_steps):
if i >= num_inference_steps - 1:
continue
t = timesteps[i]
# Expand the latents if we are doing classifier free guidance
latent_model_input = (
torch.cat([latents] * 2) if do_classifier_free_guidance else latents
)
latent_model_input = pipe.scheduler.scale_model_input(latent_model_input, t)
# Predict the noise residual
noise_pred = pipe.unet(
latent_model_input,
t,
encoder_hidden_states=text_embeddings,
added_cond_kwargs=added_cond_kwargs,
).sample
# Perform guidance
if do_classifier_free_guidance:
noise_pred_uncond, noise_pred_text = noise_pred.chunk(2)
noise_pred = noise_pred_uncond + guidance_scale * (
noise_pred_text - noise_pred_uncond
)
current_t = max(0, t.item() - (1000 // num_inference_steps)) # t
next_t = t # min(999, t.item() + (1000 // num_inference_steps)) # t+1
alpha_t = pipe.scheduler.alphas_cumprod[current_t]
alpha_t_next = pipe.scheduler.alphas_cumprod[next_t]
# Inverted update step (re-arranging the update step to get x(t) (new latents) as a function of x(t-1) (current latents)
latents = (latents - (1 - alpha_t).sqrt() * noise_pred) * (
alpha_t_next.sqrt() / alpha_t.sqrt()
) + (1 - alpha_t_next).sqrt() * noise_pred
# Store
intermediate_latents.append(latents)
return torch.cat(intermediate_latents)
def style_image_with_inversion(
pipe,
input_image,
input_image_prompt,
style_prompt,
num_steps=100,
start_step=30,
guidance_scale=3.5,
disentangle=False,
share_attn=False,
share_cross_attn=False,
share_resnet_layers=[0, 1],
share_attn_layers=[],
c2s_layers=[0, 1],
share_key=True,
share_query=True,
share_value=False,
use_adain=True,
use_content_anchor=True,
output_dir: str = None,
resnet_mode: str = None,
return_intermediate=False,
intermediate_latents=None,
):
with torch.no_grad():
pipe.vae.to(dtype=torch.float32)
latent = pipe.vae.encode(input_image.to(device) * 2 - 1)
# latent = pipe.vae.encode(input_image.to(device))
l = pipe.vae.config.scaling_factor * latent.latent_dist.sample()
if isinstance(pipe, StableDiffusionXLPipeline):
pipe.vae.to(dtype=torch.float16)
if intermediate_latents is None:
inverted_latents = invert(
pipe, l, input_image_prompt, num_inference_steps=num_steps
)
else:
inverted_latents = intermediate_latents
attn_injection.register_attention_processors(
pipe,
base_dir=output_dir,
resnet_mode=resnet_mode,
attn_mode="artist" if disentangle else "pnp",
disentangle=disentangle,
share_resblock=True,
share_attn=share_attn,
share_cross_attn=share_cross_attn,
share_resnet_layers=share_resnet_layers,
share_attn_layers=share_attn_layers,
share_key=share_key,
share_query=share_query,
share_value=share_value,
use_adain=use_adain,
c2s_layers=c2s_layers,
)
if disentangle:
final_im = sample_disentangled(
pipe,
style_prompt,
start_latents=inverted_latents[-(start_step + 1)][None],
intermediate_latents=inverted_latents,
start_step=start_step,
num_inference_steps=num_steps,
guidance_scale=guidance_scale,
use_content_anchor=use_content_anchor,
)
else:
final_im = sample(
pipe,
style_prompt,
start_latents=inverted_latents[-(start_step + 1)][None],
intermediate_latents=inverted_latents,
start_step=start_step,
num_inference_steps=num_steps,
guidance_scale=guidance_scale,
)
# unset the attention processors
attn_injection.unset_attention_processors(
pipe,
unset_share_attn=True,
unset_share_resblock=True,
)
if return_intermediate:
return final_im, inverted_latents
return final_im
if __name__ == "__main__":
# pipe = DiffusionPipeline.from_pretrained(
# # "playgroundai/playground-v2-1024px-aesthetic",
# torch_dtype=torch.float16,
# use_safetensors=True,
# add_watermarker=False,
# variant="fp16",
# )
# pipe.to("cuda")
parser = argparse.ArgumentParser(description="Stable Diffusion with OmegaConf")
parser.add_argument(
"--config", type=str, default="config.yaml", help="Path to the config file"
)
parser.add_argument(
"--mode",
type=str,
default="dataset",
choices=["dataset", "cli", "app"],
help="Path to the config file",
)
parser.add_argument(
"--image_dir", type=str, default="test.png", help="Path to the image"
)
parser.add_argument(
"--prompt",
type=str,
default="an impressionist painting",
help="Stylization prompt",
)
# mode = "single_control_content"
args = parser.parse_args()
config_dir = args.config
mode = args.mode
# mode = "dataset"
out_name = ["content_delegation", "style_delegation", "style_out"]
if mode == "app":
# gradio
import gradio as gr
import spaces
# Load a pipeline
pipe = StableDiffusionPipeline.from_pretrained(
"stabilityai/stable-diffusion-2-1-base"
).to(device)
# Set up a DDIM scheduler
pipe.scheduler = DDIMScheduler.from_config(pipe.scheduler.config)
@spaces.GPU()
def style_transfer_app(
prompt,
image,
cfg_scale=7.5,
num_content_layers=4,
num_style_layers=9,
seed=0,
progress=gr.Progress(track_tqdm=True),
):
utils.exp_utils.seed_all(seed)
image = utils.exp_utils.process_image(image, device, 512)
tgt_prompt = prompt
src_prompt = ""
prompt_in = [
"", # reconstruction
tgt_prompt, # uncontrolled style
"", # controlled style
]
share_resnet_layers = (
list(range(num_content_layers)) if num_content_layers != 0 else None
)
share_attn_layers = (
list(range(num_style_layers)) if num_style_layers != 0 else None
)
imgs = style_image_with_inversion(
pipe,
image,
src_prompt,
style_prompt=prompt_in,
num_steps=50,
start_step=0,
guidance_scale=cfg_scale,
disentangle=True,
resnet_mode="hidden",
share_attn=True,
share_cross_attn=True,
share_resnet_layers=share_resnet_layers,
share_attn_layers=share_attn_layers,
share_key=True,
share_query=True,
share_value=False,
use_content_anchor=True,
use_adain=True,
output_dir="./",
)
return imgs[2]
# load examples
examples = []
annotation = json.load(open("data/example/annotation.json"))
for entry in annotation:
image = utils.exp_utils.get_processed_image(
entry["image_path"], device, 512
)
image = transforms.ToPILImage()(image[0])
examples.append([entry["target_prompt"], image, None, None, None])
with gr.Blocks() as app:
with gr.Column():
gr.Markdown("# Artist: Aesthetically Controllable Text-Driven Stylization without Training")
gr.Markdown("## Interactive Demo, HF space version")
gr.HTML("""
<div style="display:flex;column-gap:4px;">
<a href='https://diffusionartist.github.io/'>
<img src='https://img.shields.io/badge/Project-Page-green'>
</a>
<a href='https://github.com/songrise/Artist'>
<img src='https://img.shields.io/badge/Code-github-blue'>
</a>
<a href='https://arxiv.org/abs/2407.15842'>
<img src='https://img.shields.io/badge/Paper-Arxiv-red'>
</a>
<a href='https://huggingface.co/papers/2407.15842'>
<img src='https://img.shields.io/badge/Papers-HF-ffd21f'>
</a>
</div>
""")
with gr.Row():
with gr.Column():
image_input = gr.Image(
label="Content image (will be resized to 512x512)",
interactive=True,
)
text_input = gr.Textbox(
value="An impressionist painting",
label="Text Prompt",
info="Describe the style you want to apply to the image, do not include the description of the image content itself",
lines=2,
placeholder="Enter a text prompt",
)
with gr.Accordion("Advanced settings"):
with gr.Column():
cfg_slider = gr.Slider(
0,
15,
value=7.5,
label="Classifier Free Guidance (CFG) Scale",
info="higher values give more style, 7.5 should be good for most cases",
)
content_slider = gr.Slider(
0,
9,
value=4,
step=1,
label="Number of content control layer",
info="higher values make it more similar to original image. Default to control first 4 layers",
)
style_slider = gr.Slider(
0,
9,
value=9,
step=1,
label="Number of style control layer",
info="higher values make it more similar to target style. Default to control first 9 layers, usually not necessary to change.",
)
seed_slider = gr.Slider(
0,
100,
value=0,
step=1,
label="Seed",
info="Random seed for the model",
)
submit_btn = gr.Button("Submit")
with gr.Column():
image_output= gr.Image(format="png")
gr.Examples(
examples = examples,
fn = style_transfer_app,
inputs = [text_input, image_input],
outputs = [image_output],
cache_examples=False
)
submit_btn.click(
fn=style_transfer_app,
inputs=[
text_input,
image_input,
cfg_slider,
content_slider,
style_slider,
seed_slider,
],
outputs=[image_output],
show_api=False
)
app.launch(show_api=False, show_error=True)
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