Stable Diffusion XS
Model Details
Stable Diffusion XS (SDXS) is a modified version stable diffusion for fast inference.
Usage
from diffusers import DiffusionPipeline
import torch
MODEL_PATH = "sdxs"
base = DiffusionPipeline.from_pretrained(
MODEL_PATH,
trust_remote_code=True,
torch_dtype=torch.float16,
variant="fp16",
use_safetensors=True
).to("cuda")
prompt = "柴犬、カラフルアート"
negative_prompt = ""
def tokenize_prompt(tokenizer, prompt):
text_inputs = tokenizer(
prompt,
padding="max_length",
max_length=tokenizer.model_max_length,
truncation=True,
return_tensors="pt",
)
text_input_ids = text_inputs.input_ids
return text_input_ids
def encode_prompt(text_encoders, tokenizers, prompt, hidden_size, model_max_length=77 ):
prompt_embeds_list = []
for i, text_encoder in enumerate(text_encoders):
if text_encoder is not None:
tokenizer = tokenizers[i]
text_input_ids = tokenize_prompt(tokenizer, prompt)
prompt_embeds = text_encoder(
text_input_ids.to(text_encoders[i].device), output_hidden_states=True, return_dict=False
)
pooled_prompt_embeds = prompt_embeds[0]
prompt_embeds = prompt_embeds[-1][-2]
else:
prompt_embeds = torch.zeros((1, model_max_length, hidden_size))
pooled_prompt_embeds = torch.zeros((1, hidden_size))
# We are only ALWAYS interested in the pooled output of the final text encoder
prompt_embeds = prompt_embeds.to("cuda")
bs_embed, seq_len, _ = prompt_embeds.shape
prompt_embeds = prompt_embeds.view(bs_embed, seq_len, -1)
prompt_embeds_list.append(prompt_embeds)
prompt_embeds = torch.concat(prompt_embeds_list, dim=-1)
pooled_prompt_embeds = pooled_prompt_embeds.view(bs_embed, -1)
return prompt_embeds, pooled_prompt_embeds
prompt_embeds, pooled_prompt_embeds = encode_prompt([None, base.text_encoder],[None, base.tokenizer], prompt, 768)
negative_prompt_embeds, negative_pooled_prompt_embeds = encode_prompt([None, base.text_encoder],[None, base.tokenizer], negative_prompt, 768)
#generator = [torch.Generator(device="cuda").manual_seed(i) for i in range(4)]
image = base(
prompt_embeds=prompt_embeds,
pooled_prompt_embeds=pooled_prompt_embeds,
negative_prompt_embeds=negative_prompt_embeds,
negative_pooled_prompt_embeds=negative_pooled_prompt_embeds,
num_inference_steps=20,
).images[0]
display(image)
Model Details
- Developed by: AiArtLab
- Model type: Diffusion-based text-to-image generative model
- Model Description: This model is a fine-tuned model based on colorfulxl_v27.
- License:
Uses
Direct Use
Research: possible research areas/tasks include:
- Generation of artworks and use in design and other artistic processes.
- Applications in educational or creative tools.
- Research on generative models.
- Safe deployment of models which have the potential to generate harmful content.
- Probing and understanding the limitations and biases of generative models.
Excluded uses are described below.
Out-of-Scope Use
The model was not trained to be factual or true representations of people or events, and therefore using the model to generate such content is out-of-scope for the abilities of this model.
Limitations and Bias
Limitations
- The model does not achieve perfect photorealism
- The model cannot render legible text
- The model struggles with more difficult tasks which involve compositionality, such as rendering an image corresponding to “A red cube on top of a blue sphere”
- Faces and people in general may not be generated properly.
- The autoencoding part of the model is lossy.
Bias
While the capabilities of image generation models are impressive, they can also reinforce or exacerbate social biases.
How to cite
@misc{SDXS,
url = {[https://huggingface.co/recoilme/sdxs](https://huggingface.co/recoilme/sdxs)},
title = {Stable Diffusion XS},
author = {recoilme}
}
Contact
- For questions and comments about the model, please join https://aiartlab.org/.
- For future announcements / information about AiArtLab AI models, research, and events, please follow Discord.
- For business and partnership inquiries, please contact https://t.me/recoilme
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