DPO LoRA Stable Diffusion v2-1

Model trained with LoRA implementation of Diffusion DPO Read more here

Base Model: https://huggingface.co/stabilityai/stable-diffusion-2-1

Running with 🧨 diffusers library


from diffusers import DiffusionPipeline
from diffusers.utils import make_image_grid
import torch

pipe = DiffusionPipeline.from_pretrained(
    "stabilityai/sd-turbo",  # SD Turbo is a destilled version of Stable Diffusion 2.1
    # "stabilityai/stable-diffusion-2-1", # for the original stable diffusion 2.1 model
    torch_dtype=torch.float16, variant="fp16"
)
pipe.to("cuda")
pipe.load_lora_weights("radames/sd-21-DPO-LoRA", adapter_name="dpo-lora-sd21")
pipe.set_adapters(["dpo-lora-sd21"], adapter_weights=[1.0]) # you can play with adapter_weights to increase the effect of the LoRA model
seed = 123123
prompt = "portrait headshot professional of elon musk"
negative_prompt = "3d render, cartoon, drawing, art, low light"
generator = torch.Generator().manual_seed(seed)
images = pipe(
    prompt=prompt,
    negative_prompt=negative_prompt,
    width=512,
    height=512,
    num_inference_steps=2,
    generator=generator,
    guidance_scale=1.0,
    num_images_per_prompt=4
).images
make_image_grid(images, 1, 4)