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on
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A newer version of the Gradio SDK is available:
5.4.0
DreamBooth by colossalai
DreamBooth is a method to personalize text2image models like stable diffusion given just a few(3~5) images of a subject.
The train_dreambooth_colossalai.py
script shows how to implement the training procedure and adapt it for stable diffusion.
By accommodating model data in CPU and GPU and moving the data to the computing device when necessary, Gemini, the Heterogeneous Memory Manager of Colossal-AI can breakthrough the GPU memory wall by using GPU and CPU memory (composed of CPU DRAM or nvme SSD memory) together at the same time. Moreover, the model scale can be further improved by combining heterogeneous training with the other parallel approaches, such as data parallel, tensor parallel and pipeline parallel.
Installing the dependencies
Before running the scripts, make sure to install the library's training dependencies:
pip install -r requirements.txt
Install ColossalAI
From PyPI
pip install colossalai
From source
git clone https://github.com/hpcaitech/ColossalAI.git
cd ColossalAI
# install colossalai
pip install .
Dataset for Teyvat BLIP captions
Dataset used to train Teyvat characters text to image model.
BLIP generated captions for characters images from genshin-impact fandom wikiand biligame wiki for genshin impact.
For each row the dataset contains image
and text
keys. image
is a varying size PIL png, and text
is the accompanying text caption. Only a train split is provided.
The text
include the tag Teyvat
, Name
,Element
, Weapon
, Region
, Model type
, and Description
, the Description
is captioned with the pre-trained BLIP model.
Training
The arguement placement
can be cpu
, auto
, cuda
, with cpu
the GPU RAM required can be minimized to 4GB but will deceleration, with cuda
you can also reduce GPU memory by half but accelerated training, with auto
a more balanced solution for speed and memory can be obtained。
Note: Change the resolution
to 768 if you are using the stable-diffusion-2 768x768 model.
export MODEL_NAME="CompVis/stable-diffusion-v1-4"
export INSTANCE_DIR="path-to-instance-images"
export OUTPUT_DIR="path-to-save-model"
torchrun --nproc_per_node 2 train_dreambooth_colossalai.py \
--pretrained_model_name_or_path=$MODEL_NAME \
--instance_data_dir=$INSTANCE_DIR \
--output_dir=$OUTPUT_DIR \
--instance_prompt="a photo of sks dog" \
--resolution=512 \
--train_batch_size=1 \
--learning_rate=5e-6 \
--lr_scheduler="constant" \
--lr_warmup_steps=0 \
--max_train_steps=400 \
--placement="cuda"
Training with prior-preservation loss
Prior-preservation is used to avoid overfitting and language-drift. Refer to the paper to learn more about it. For prior-preservation we first generate images using the model with a class prompt and then use those during training along with our data.
According to the paper, it's recommended to generate num_epochs * num_samples
images for prior-preservation. 200-300 works well for most cases. The num_class_images
flag sets the number of images to generate with the class prompt. You can place existing images in class_data_dir
, and the training script will generate any additional images so that num_class_images
are present in class_data_dir
during training time.
export MODEL_NAME="CompVis/stable-diffusion-v1-4"
export INSTANCE_DIR="path-to-instance-images"
export CLASS_DIR="path-to-class-images"
export OUTPUT_DIR="path-to-save-model"
torchrun --nproc_per_node 2 train_dreambooth_colossalai.py \
--pretrained_model_name_or_path=$MODEL_NAME \
--instance_data_dir=$INSTANCE_DIR \
--class_data_dir=$CLASS_DIR \
--output_dir=$OUTPUT_DIR \
--with_prior_preservation --prior_loss_weight=1.0 \
--instance_prompt="a photo of sks dog" \
--class_prompt="a photo of dog" \
--resolution=512 \
--train_batch_size=1 \
--learning_rate=5e-6 \
--lr_scheduler="constant" \
--lr_warmup_steps=0 \
--max_train_steps=800 \
--placement="cuda"
Inference
Once you have trained a model using above command, the inference can be done simply using the StableDiffusionPipeline
. Make sure to include the identifier
(e.g. sks in above example) in your prompt.
from diffusers import StableDiffusionPipeline
import torch
model_id = "path-to-save-model"
pipe = StableDiffusionPipeline.from_pretrained(model_id, torch_dtype=torch.float16).to("cuda")
prompt = "A photo of sks dog in a bucket"
image = pipe(prompt, num_inference_steps=50, guidance_scale=7.5).images[0]
image.save("dog-bucket.png")