ControlNet with Stable Diffusion XL
ControlNet was introduced in Adding Conditional Control to Text-to-Image Diffusion Models by Lvmin Zhang, Anyi Rao, and Maneesh Agrawala.
With a ControlNet model, you can provide an additional control image to condition and control Stable Diffusion generation. For example, if you provide a depth map, the ControlNet model generates an image that’ll preserve the spatial information from the depth map. It is a more flexible and accurate way to control the image generation process.
The abstract from the paper is:
We present ControlNet, a neural network architecture to add spatial conditioning controls to large, pretrained text-to-image diffusion models. ControlNet locks the production-ready large diffusion models, and reuses their deep and robust encoding layers pretrained with billions of images as a strong backbone to learn a diverse set of conditional controls. The neural architecture is connected with “zero convolutions” (zero-initialized convolution layers) that progressively grow the parameters from zero and ensure that no harmful noise could affect the finetuning. We test various conditioning controls, eg, edges, depth, segmentation, human pose, etc, with Stable Diffusion, using single or multiple conditions, with or without prompts. We show that the training of ControlNets is robust with small (<50k) and large (>1m) datasets. Extensive results show that ControlNet may facilitate wider applications to control image diffusion models.
You can find additional smaller Stable Diffusion XL (SDXL) ControlNet checkpoints from the 🤗 Diffusers Hub organization, and browse community-trained checkpoints on the Hub.
🧪 Many of the SDXL ControlNet checkpoints are experimental, and there is a lot of room for improvement. Feel free to open an Issue and leave us feedback on how we can improve!
If you don’t see a checkpoint you’re interested in, you can train your own SDXL ControlNet with our training script.
Make sure to check out the Schedulers guide to learn how to explore the tradeoff between scheduler speed and quality, and see the reuse components across pipelines section to learn how to efficiently load the same components into multiple pipelines.
StableDiffusionXLControlNetPipeline
class diffusers.StableDiffusionXLControlNetPipeline
< source >( vae: AutoencoderKL text_encoder: CLIPTextModel text_encoder_2: CLIPTextModelWithProjection tokenizer: CLIPTokenizer tokenizer_2: CLIPTokenizer unet: UNet2DConditionModel controlnet: Union scheduler: KarrasDiffusionSchedulers force_zeros_for_empty_prompt: bool = True add_watermarker: Optional = None feature_extractor: CLIPImageProcessor = None image_encoder: CLIPVisionModelWithProjection = None )
Parameters
- vae (AutoencoderKL) — Variational Auto-Encoder (VAE) model to encode and decode images to and from latent representations.
- text_encoder (CLIPTextModel) — Frozen text-encoder (clip-vit-large-patch14).
- text_encoder_2 (CLIPTextModelWithProjection) — Second frozen text-encoder (laion/CLIP-ViT-bigG-14-laion2B-39B-b160k).
- tokenizer (CLIPTokenizer) —
A
CLIPTokenizer
to tokenize text. - tokenizer_2 (CLIPTokenizer) —
A
CLIPTokenizer
to tokenize text. - unet (UNet2DConditionModel) —
A
UNet2DConditionModel
to denoise the encoded image latents. - controlnet (ControlNetModel or
List[ControlNetModel]
) — Provides additional conditioning to theunet
during the denoising process. If you set multiple ControlNets as a list, the outputs from each ControlNet are added together to create one combined additional conditioning. - scheduler (SchedulerMixin) —
A scheduler to be used in combination with
unet
to denoise the encoded image latents. Can be one of DDIMScheduler, LMSDiscreteScheduler, or PNDMScheduler. - force_zeros_for_empty_prompt (
bool
, optional, defaults to"True"
) — Whether the negative prompt embeddings should always be set to 0. Also see the config ofstabilityai/stable-diffusion-xl-base-1-0
. - add_watermarker (
bool
, optional) — Whether to use the invisible_watermark library to watermark output images. If not defined, it defaults toTrue
if the package is installed; otherwise no watermarker is used.
Pipeline for text-to-image generation using Stable Diffusion XL with ControlNet guidance.
This model inherits from DiffusionPipeline. Check the superclass documentation for the generic methods implemented for all pipelines (downloading, saving, running on a particular device, etc.).
The pipeline also inherits the following loading methods:
- load_textual_inversion() for loading textual inversion embeddings
- load_lora_weights() for loading LoRA weights
save_lora_weights()
for saving LoRA weights- from_single_file() for loading
.ckpt
files - load_ip_adapter() for loading IP Adapters
__call__
< source >( prompt: Union = None prompt_2: Union = None image: Union = None height: Optional = None width: Optional = None num_inference_steps: int = 50 guidance_scale: float = 5.0 negative_prompt: Union = None negative_prompt_2: Union = None num_images_per_prompt: Optional = 1 eta: float = 0.0 generator: Union = None latents: Optional = None prompt_embeds: Optional = None negative_prompt_embeds: Optional = None pooled_prompt_embeds: Optional = None negative_pooled_prompt_embeds: Optional = None ip_adapter_image: Union = None output_type: Optional = 'pil' return_dict: bool = True cross_attention_kwargs: Optional = None controlnet_conditioning_scale: Union = 1.0 guess_mode: bool = False control_guidance_start: Union = 0.0 control_guidance_end: Union = 1.0 original_size: Tuple = None crops_coords_top_left: Tuple = (0, 0) target_size: Tuple = None negative_original_size: Optional = None negative_crops_coords_top_left: Tuple = (0, 0) negative_target_size: Optional = None clip_skip: Optional = None callback_on_step_end: Optional = None callback_on_step_end_tensor_inputs: List = ['latents'] **kwargs ) → StableDiffusionPipelineOutput or tuple
Parameters
- prompt (
str
orList[str]
, optional) — The prompt or prompts to guide image generation. If not defined, you need to passprompt_embeds
. - prompt_2 (
str
orList[str]
, optional) — The prompt or prompts to be sent totokenizer_2
andtext_encoder_2
. If not defined,prompt
is used in both text-encoders. - image (
torch.FloatTensor
,PIL.Image.Image
,np.ndarray
,List[torch.FloatTensor]
,List[PIL.Image.Image]
,List[np.ndarray]
, —List[List[torch.FloatTensor]]
,List[List[np.ndarray]]
orList[List[PIL.Image.Image]]
): The ControlNet input condition to provide guidance to theunet
for generation. If the type is specified astorch.FloatTensor
, it is passed to ControlNet as is.PIL.Image.Image
can also be accepted as an image. The dimensions of the output image defaults toimage
’s dimensions. If height and/or width are passed,image
is resized accordingly. If multiple ControlNets are specified ininit
, images must be passed as a list such that each element of the list can be correctly batched for input to a single ControlNet. - height (
int
, optional, defaults toself.unet.config.sample_size * self.vae_scale_factor
) — The height in pixels of the generated image. Anything below 512 pixels won’t work well for stabilityai/stable-diffusion-xl-base-1.0 and checkpoints that are not specifically fine-tuned on low resolutions. - width (
int
, optional, defaults toself.unet.config.sample_size * self.vae_scale_factor
) — The width in pixels of the generated image. Anything below 512 pixels won’t work well for stabilityai/stable-diffusion-xl-base-1.0 and checkpoints that are not specifically fine-tuned on low resolutions. - num_inference_steps (
int
, optional, defaults to 50) — The number of denoising steps. More denoising steps usually lead to a higher quality image at the expense of slower inference. - guidance_scale (
float
, optional, defaults to 5.0) — A higher guidance scale value encourages the model to generate images closely linked to the textprompt
at the expense of lower image quality. Guidance scale is enabled whenguidance_scale > 1
. - negative_prompt (
str
orList[str]
, optional) — The prompt or prompts to guide what to not include in image generation. If not defined, you need to passnegative_prompt_embeds
instead. Ignored when not using guidance (guidance_scale < 1
). - negative_prompt_2 (
str
orList[str]
, optional) — The prompt or prompts to guide what to not include in image generation. This is sent totokenizer_2
andtext_encoder_2
. If not defined,negative_prompt
is used in both text-encoders. - num_images_per_prompt (
int
, optional, defaults to 1) — The number of images to generate per prompt. - eta (
float
, optional, defaults to 0.0) — Corresponds to parameter eta (η) from the DDIM paper. Only applies to the DDIMScheduler, and is ignored in other schedulers. - generator (
torch.Generator
orList[torch.Generator]
, optional) — Atorch.Generator
to make generation deterministic. - latents (
torch.FloatTensor
, optional) — Pre-generated noisy latents sampled from a Gaussian distribution, to be used as inputs for image generation. Can be used to tweak the same generation with different prompts. If not provided, a latents tensor is generated by sampling using the supplied randomgenerator
. - prompt_embeds (
torch.FloatTensor
, optional) — Pre-generated text embeddings. Can be used to easily tweak text inputs (prompt weighting). If not provided, text embeddings are generated from theprompt
input argument. - negative_prompt_embeds (
torch.FloatTensor
, optional) — Pre-generated negative text embeddings. Can be used to easily tweak text inputs (prompt weighting). If not provided,negative_prompt_embeds
are generated from thenegative_prompt
input argument. - pooled_prompt_embeds (
torch.FloatTensor
, optional) — Pre-generated pooled text embeddings. Can be used to easily tweak text inputs (prompt weighting). If not provided, pooled text embeddings are generated fromprompt
input argument. - negative_pooled_prompt_embeds (
torch.FloatTensor
, optional) — Pre-generated negative pooled text embeddings. Can be used to easily tweak text inputs (prompt weighting). If not provided, poolednegative_prompt_embeds
are generated fromnegative_prompt
input argument. ip_adapter_image — (PipelineImageInput
, optional): Optional image input to work with IP Adapters. - output_type (
str
, optional, defaults to"pil"
) — The output format of the generated image. Choose betweenPIL.Image
ornp.array
. - return_dict (
bool
, optional, defaults toTrue
) — Whether or not to return a StableDiffusionPipelineOutput instead of a plain tuple. - cross_attention_kwargs (
dict
, optional) — A kwargs dictionary that if specified is passed along to theAttentionProcessor
as defined inself.processor
. - controlnet_conditioning_scale (
float
orList[float]
, optional, defaults to 1.0) — The outputs of the ControlNet are multiplied bycontrolnet_conditioning_scale
before they are added to the residual in the originalunet
. If multiple ControlNets are specified ininit
, you can set the corresponding scale as a list. - guess_mode (
bool
, optional, defaults toFalse
) — The ControlNet encoder tries to recognize the content of the input image even if you remove all prompts. Aguidance_scale
value between 3.0 and 5.0 is recommended. - control_guidance_start (
float
orList[float]
, optional, defaults to 0.0) — The percentage of total steps at which the ControlNet starts applying. - control_guidance_end (
float
orList[float]
, optional, defaults to 1.0) — The percentage of total steps at which the ControlNet stops applying. - original_size (
Tuple[int]
, optional, defaults to (1024, 1024)) — Iforiginal_size
is not the same astarget_size
the image will appear to be down- or upsampled.original_size
defaults to(height, width)
if not specified. Part of SDXL’s micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952. - crops_coords_top_left (
Tuple[int]
, optional, defaults to (0, 0)) —crops_coords_top_left
can be used to generate an image that appears to be “cropped” from the positioncrops_coords_top_left
downwards. Favorable, well-centered images are usually achieved by settingcrops_coords_top_left
to (0, 0). Part of SDXL’s micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952. - target_size (
Tuple[int]
, optional, defaults to (1024, 1024)) — For most cases,target_size
should be set to the desired height and width of the generated image. If not specified it will default to(height, width)
. Part of SDXL’s micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952. - negative_original_size (
Tuple[int]
, optional, defaults to (1024, 1024)) — To negatively condition the generation process based on a specific image resolution. Part of SDXL’s micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952. For more information, refer to this issue thread: https://github.com/huggingface/diffusers/issues/4208. - negative_crops_coords_top_left (
Tuple[int]
, optional, defaults to (0, 0)) — To negatively condition the generation process based on a specific crop coordinates. Part of SDXL’s micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952. For more information, refer to this issue thread: https://github.com/huggingface/diffusers/issues/4208. - negative_target_size (
Tuple[int]
, optional, defaults to (1024, 1024)) — To negatively condition the generation process based on a target image resolution. It should be as same as thetarget_size
for most cases. Part of SDXL’s micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952. For more information, refer to this issue thread: https://github.com/huggingface/diffusers/issues/4208. - clip_skip (
int
, optional) — Number of layers to be skipped from CLIP while computing the prompt embeddings. A value of 1 means that the output of the pre-final layer will be used for computing the prompt embeddings. - callback_on_step_end (
Callable
, optional) — A function that calls at the end of each denoising steps during the inference. The function is called with the following arguments:callback_on_step_end(self: DiffusionPipeline, step: int, timestep: int, callback_kwargs: Dict)
.callback_kwargs
will include a list of all tensors as specified bycallback_on_step_end_tensor_inputs
. - callback_on_step_end_tensor_inputs (
List
, optional) — The list of tensor inputs for thecallback_on_step_end
function. The tensors specified in the list will be passed ascallback_kwargs
argument. You will only be able to include variables listed in the._callback_tensor_inputs
attribute of your pipeine class.
Returns
StableDiffusionPipelineOutput or tuple
If return_dict
is True
, StableDiffusionPipelineOutput is returned,
otherwise a tuple
is returned containing the output images.
The call function to the pipeline for generation.
Examples:
>>> # !pip install opencv-python transformers accelerate
>>> from diffusers import StableDiffusionXLControlNetPipeline, ControlNetModel, AutoencoderKL
>>> from diffusers.utils import load_image
>>> import numpy as np
>>> import torch
>>> import cv2
>>> from PIL import Image
>>> prompt = "aerial view, a futuristic research complex in a bright foggy jungle, hard lighting"
>>> negative_prompt = "low quality, bad quality, sketches"
>>> # download an image
>>> image = load_image(
... "https://hf.co/datasets/hf-internal-testing/diffusers-images/resolve/main/sd_controlnet/hf-logo.png"
... )
>>> # initialize the models and pipeline
>>> controlnet_conditioning_scale = 0.5 # recommended for good generalization
>>> controlnet = ControlNetModel.from_pretrained(
... "diffusers/controlnet-canny-sdxl-1.0", torch_dtype=torch.float16
... )
>>> vae = AutoencoderKL.from_pretrained("madebyollin/sdxl-vae-fp16-fix", torch_dtype=torch.float16)
>>> pipe = StableDiffusionXLControlNetPipeline.from_pretrained(
... "stabilityai/stable-diffusion-xl-base-1.0", controlnet=controlnet, vae=vae, torch_dtype=torch.float16
... )
>>> pipe.enable_model_cpu_offload()
>>> # get canny image
>>> image = np.array(image)
>>> image = cv2.Canny(image, 100, 200)
>>> image = image[:, :, None]
>>> image = np.concatenate([image, image, image], axis=2)
>>> canny_image = Image.fromarray(image)
>>> # generate image
>>> image = pipe(
... prompt, controlnet_conditioning_scale=controlnet_conditioning_scale, image=canny_image
... ).images[0]
Disables the FreeU mechanism if enabled.
Disable sliced VAE decoding. If enable_vae_slicing
was previously enabled, this method will go back to
computing decoding in one step.
Disable tiled VAE decoding. If enable_vae_tiling
was previously enabled, this method will go back to
computing decoding in one step.
enable_freeu
< source >( s1: float s2: float b1: float b2: float )
Parameters
- s1 (
float
) — Scaling factor for stage 1 to attenuate the contributions of the skip features. This is done to mitigate “oversmoothing effect” in the enhanced denoising process. - s2 (
float
) — Scaling factor for stage 2 to attenuate the contributions of the skip features. This is done to mitigate “oversmoothing effect” in the enhanced denoising process. - b1 (
float
) — Scaling factor for stage 1 to amplify the contributions of backbone features. - b2 (
float
) — Scaling factor for stage 2 to amplify the contributions of backbone features.
Enables the FreeU mechanism as in https://arxiv.org/abs/2309.11497.
The suffixes after the scaling factors represent the stages where they are being applied.
Please refer to the official repository for combinations of the values that are known to work well for different pipelines such as Stable Diffusion v1, v2, and Stable Diffusion XL.
Enable sliced VAE decoding. When this option is enabled, the VAE will split the input tensor in slices to compute decoding in several steps. This is useful to save some memory and allow larger batch sizes.
Enable tiled VAE decoding. When this option is enabled, the VAE will split the input tensor into tiles to compute decoding and encoding in several steps. This is useful for saving a large amount of memory and to allow processing larger images.
encode_prompt
< source >( prompt: str prompt_2: Optional = None device: Optional = None num_images_per_prompt: int = 1 do_classifier_free_guidance: bool = True negative_prompt: Optional = None negative_prompt_2: Optional = None prompt_embeds: Optional = None negative_prompt_embeds: Optional = None pooled_prompt_embeds: Optional = None negative_pooled_prompt_embeds: Optional = None lora_scale: Optional = None clip_skip: Optional = None )
Parameters
- prompt (
str
orList[str]
, optional) — prompt to be encoded - prompt_2 (
str
orList[str]
, optional) — The prompt or prompts to be sent to thetokenizer_2
andtext_encoder_2
. If not defined,prompt
is used in both text-encoders device — (torch.device
): torch device - num_images_per_prompt (
int
) — number of images that should be generated per prompt - do_classifier_free_guidance (
bool
) — whether to use classifier free guidance or not - negative_prompt (
str
orList[str]
, optional) — The prompt or prompts not to guide the image generation. If not defined, one has to passnegative_prompt_embeds
instead. Ignored when not using guidance (i.e., ignored ifguidance_scale
is less than1
). - negative_prompt_2 (
str
orList[str]
, optional) — The prompt or prompts not to guide the image generation to be sent totokenizer_2
andtext_encoder_2
. If not defined,negative_prompt
is used in both text-encoders - prompt_embeds (
torch.FloatTensor
, optional) — Pre-generated text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not provided, text embeddings will be generated fromprompt
input argument. - negative_prompt_embeds (
torch.FloatTensor
, optional) — Pre-generated negative text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not provided, negative_prompt_embeds will be generated fromnegative_prompt
input argument. - pooled_prompt_embeds (
torch.FloatTensor
, optional) — Pre-generated pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not provided, pooled text embeddings will be generated fromprompt
input argument. - negative_pooled_prompt_embeds (
torch.FloatTensor
, optional) — Pre-generated negative pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not provided, pooled negative_prompt_embeds will be generated fromnegative_prompt
input argument. - lora_scale (
float
, optional) — A lora scale that will be applied to all LoRA layers of the text encoder if LoRA layers are loaded. - clip_skip (
int
, optional) — Number of layers to be skipped from CLIP while computing the prompt embeddings. A value of 1 means that the output of the pre-final layer will be used for computing the prompt embeddings.
Encodes the prompt into text encoder hidden states.
get_guidance_scale_embedding
< source >( w embedding_dim = 512 dtype = torch.float32 ) → torch.FloatTensor
Parameters
- timesteps (
torch.Tensor
) — generate embedding vectors at these timesteps - embedding_dim (
int
, optional, defaults to 512) — dimension of the embeddings to generate dtype — data type of the generated embeddings
Returns
torch.FloatTensor
Embedding vectors with shape (len(timesteps), embedding_dim)
StableDiffusionXLControlNetImg2ImgPipeline
class diffusers.StableDiffusionXLControlNetImg2ImgPipeline
< source >( vae: AutoencoderKL text_encoder: CLIPTextModel text_encoder_2: CLIPTextModelWithProjection tokenizer: CLIPTokenizer tokenizer_2: CLIPTokenizer unet: UNet2DConditionModel controlnet: Union scheduler: KarrasDiffusionSchedulers requires_aesthetics_score: bool = False force_zeros_for_empty_prompt: bool = True add_watermarker: Optional = None feature_extractor: CLIPImageProcessor = None image_encoder: CLIPVisionModelWithProjection = None )
Parameters
- vae (AutoencoderKL) — Variational Auto-Encoder (VAE) Model to encode and decode images to and from latent representations.
- text_encoder (
CLIPTextModel
) — Frozen text-encoder. Stable Diffusion uses the text portion of CLIP, specifically the clip-vit-large-patch14 variant. - text_encoder_2 (
CLIPTextModelWithProjection
) — Second frozen text-encoder. Stable Diffusion XL uses the text and pool portion of CLIP, specifically the laion/CLIP-ViT-bigG-14-laion2B-39B-b160k variant. - tokenizer (
CLIPTokenizer
) — Tokenizer of class CLIPTokenizer. - tokenizer_2 (
CLIPTokenizer
) — Second Tokenizer of class CLIPTokenizer. - unet (UNet2DConditionModel) — Conditional U-Net architecture to denoise the encoded image latents.
- controlnet (ControlNetModel or
List[ControlNetModel]
) — Provides additional conditioning to the unet during the denoising process. If you set multiple ControlNets as a list, the outputs from each ControlNet are added together to create one combined additional conditioning. - scheduler (SchedulerMixin) —
A scheduler to be used in combination with
unet
to denoise the encoded image latents. Can be one of DDIMScheduler, LMSDiscreteScheduler, or PNDMScheduler. - requires_aesthetics_score (
bool
, optional, defaults to"False"
) — Whether theunet
requires anaesthetic_score
condition to be passed during inference. Also see the config ofstabilityai/stable-diffusion-xl-refiner-1-0
. - force_zeros_for_empty_prompt (
bool
, optional, defaults to"True"
) — Whether the negative prompt embeddings shall be forced to always be set to 0. Also see the config ofstabilityai/stable-diffusion-xl-base-1-0
. - add_watermarker (
bool
, optional) — Whether to use the invisible_watermark library to watermark output images. If not defined, it will default to True if the package is installed, otherwise no watermarker will be used. - feature_extractor (CLIPImageProcessor) —
A
CLIPImageProcessor
to extract features from generated images; used as inputs to thesafety_checker
.
Pipeline for image-to-image generation using Stable Diffusion XL with ControlNet guidance.
This model inherits from DiffusionPipeline. Check the superclass documentation for the generic methods the library implements for all the pipelines (such as downloading or saving, running on a particular device, etc.)
The pipeline also inherits the following loading methods:
- load_textual_inversion() for loading textual inversion embeddings
- load_lora_weights() for loading LoRA weights
save_lora_weights()
for saving LoRA weights- load_ip_adapter() for loading IP Adapters
__call__
< source >( prompt: Union = None prompt_2: Union = None image: Union = None control_image: Union = None height: Optional = None width: Optional = None strength: float = 0.8 num_inference_steps: int = 50 guidance_scale: float = 5.0 negative_prompt: Union = None negative_prompt_2: Union = None num_images_per_prompt: Optional = 1 eta: float = 0.0 generator: Union = None latents: Optional = None prompt_embeds: Optional = None negative_prompt_embeds: Optional = None pooled_prompt_embeds: Optional = None negative_pooled_prompt_embeds: Optional = None ip_adapter_image: Union = None output_type: Optional = 'pil' return_dict: bool = True cross_attention_kwargs: Optional = None controlnet_conditioning_scale: Union = 0.8 guess_mode: bool = False control_guidance_start: Union = 0.0 control_guidance_end: Union = 1.0 original_size: Tuple = None crops_coords_top_left: Tuple = (0, 0) target_size: Tuple = None negative_original_size: Optional = None negative_crops_coords_top_left: Tuple = (0, 0) negative_target_size: Optional = None aesthetic_score: float = 6.0 negative_aesthetic_score: float = 2.5 clip_skip: Optional = None callback_on_step_end: Optional = None callback_on_step_end_tensor_inputs: List = ['latents'] **kwargs ) → StableDiffusionPipelineOutput or tuple
Parameters
- prompt (
str
orList[str]
, optional) — The prompt or prompts to guide the image generation. If not defined, one has to passprompt_embeds
. instead. - prompt_2 (
str
orList[str]
, optional) — The prompt or prompts to be sent to thetokenizer_2
andtext_encoder_2
. If not defined,prompt
is used in both text-encoders - image (
torch.FloatTensor
,PIL.Image.Image
,np.ndarray
,List[torch.FloatTensor]
,List[PIL.Image.Image]
,List[np.ndarray]
, —List[List[torch.FloatTensor]]
,List[List[np.ndarray]]
orList[List[PIL.Image.Image]]
): The initial image will be used as the starting point for the image generation process. Can also accept image latents asimage
, if passing latents directly, it will not be encoded again. - control_image (
torch.FloatTensor
,PIL.Image.Image
,np.ndarray
,List[torch.FloatTensor]
,List[PIL.Image.Image]
,List[np.ndarray]
, —List[List[torch.FloatTensor]]
,List[List[np.ndarray]]
orList[List[PIL.Image.Image]]
): The ControlNet input condition. ControlNet uses this input condition to generate guidance to Unet. If the type is specified asTorch.FloatTensor
, it is passed to ControlNet as is.PIL.Image.Image
can also be accepted as an image. The dimensions of the output image defaults toimage
’s dimensions. If height and/or width are passed,image
is resized according to them. If multiple ControlNets are specified in init, images must be passed as a list such that each element of the list can be correctly batched for input to a single controlnet. - height (
int
, optional, defaults to the size of control_image) — The height in pixels of the generated image. Anything below 512 pixels won’t work well for stabilityai/stable-diffusion-xl-base-1.0 and checkpoints that are not specifically fine-tuned on low resolutions. - width (
int
, optional, defaults to the size of control_image) — The width in pixels of the generated image. Anything below 512 pixels won’t work well for stabilityai/stable-diffusion-xl-base-1.0 and checkpoints that are not specifically fine-tuned on low resolutions. - num_inference_steps (
int
, optional, defaults to 50) — The number of denoising steps. More denoising steps usually lead to a higher quality image at the expense of slower inference. - strength (
float
, optional, defaults to 0.3) — Conceptually, indicates how much to transform the referenceimage
. Must be between 0 and 1.image
will be used as a starting point, adding more noise to it the larger thestrength
. The number of denoising steps depends on the amount of noise initially added. Whenstrength
is 1, added noise will be maximum and the denoising process will run for the full number of iterations specified innum_inference_steps
. - guidance_scale (
float
, optional, defaults to 7.5) — Guidance scale as defined in Classifier-Free Diffusion Guidance.guidance_scale
is defined asw
of equation 2. of Imagen Paper. Guidance scale is enabled by settingguidance_scale > 1
. Higher guidance scale encourages to generate images that are closely linked to the textprompt
, usually at the expense of lower image quality. - negative_prompt (
str
orList[str]
, optional) — The prompt or prompts not to guide the image generation. If not defined, one has to passnegative_prompt_embeds
instead. Ignored when not using guidance (i.e., ignored ifguidance_scale
is less than1
). - negative_prompt_2 (
str
orList[str]
, optional) — The prompt or prompts not to guide the image generation to be sent totokenizer_2
andtext_encoder_2
. If not defined,negative_prompt
is used in both text-encoders - num_images_per_prompt (
int
, optional, defaults to 1) — The number of images to generate per prompt. - eta (
float
, optional, defaults to 0.0) — Corresponds to parameter eta (η) in the DDIM paper: https://arxiv.org/abs/2010.02502. Only applies to schedulers.DDIMScheduler, will be ignored for others. - generator (
torch.Generator
orList[torch.Generator]
, optional) — One or a list of torch generator(s) to make generation deterministic. - latents (
torch.FloatTensor
, optional) — Pre-generated noisy latents, sampled from a Gaussian distribution, to be used as inputs for image generation. Can be used to tweak the same generation with different prompts. If not provided, a latents tensor will ge generated by sampling using the supplied randomgenerator
. - prompt_embeds (
torch.FloatTensor
, optional) — Pre-generated text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not provided, text embeddings will be generated fromprompt
input argument. - negative_prompt_embeds (
torch.FloatTensor
, optional) — Pre-generated negative text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not provided, negative_prompt_embeds will be generated fromnegative_prompt
input argument. - pooled_prompt_embeds (
torch.FloatTensor
, optional) — Pre-generated pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not provided, pooled text embeddings will be generated fromprompt
input argument. - negative_pooled_prompt_embeds (
torch.FloatTensor
, optional) — Pre-generated negative pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not provided, pooled negative_prompt_embeds will be generated fromnegative_prompt
input argument. ip_adapter_image — (PipelineImageInput
, optional): Optional image input to work with IP Adapters. - output_type (
str
, optional, defaults to"pil"
) — The output format of the generate image. Choose between PIL:PIL.Image.Image
ornp.array
. - return_dict (
bool
, optional, defaults toTrue
) — Whether or not to return a StableDiffusionPipelineOutput instead of a plain tuple. - cross_attention_kwargs (
dict
, optional) — A kwargs dictionary that if specified is passed along to theAttentionProcessor
as defined underself.processor
in diffusers.models.attention_processor. - controlnet_conditioning_scale (
float
orList[float]
, optional, defaults to 1.0) — The outputs of the controlnet are multiplied bycontrolnet_conditioning_scale
before they are added to the residual in the original unet. If multiple ControlNets are specified in init, you can set the corresponding scale as a list. - guess_mode (
bool
, optional, defaults toFalse
) — In this mode, the ControlNet encoder will try best to recognize the content of the input image even if you remove all prompts. Theguidance_scale
between 3.0 and 5.0 is recommended. - control_guidance_start (
float
orList[float]
, optional, defaults to 0.0) — The percentage of total steps at which the controlnet starts applying. - control_guidance_end (
float
orList[float]
, optional, defaults to 1.0) — The percentage of total steps at which the controlnet stops applying. - original_size (
Tuple[int]
, optional, defaults to (1024, 1024)) — Iforiginal_size
is not the same astarget_size
the image will appear to be down- or upsampled.original_size
defaults to(height, width)
if not specified. Part of SDXL’s micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952. - crops_coords_top_left (
Tuple[int]
, optional, defaults to (0, 0)) —crops_coords_top_left
can be used to generate an image that appears to be “cropped” from the positioncrops_coords_top_left
downwards. Favorable, well-centered images are usually achieved by settingcrops_coords_top_left
to (0, 0). Part of SDXL’s micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952. - target_size (
Tuple[int]
, optional, defaults to (1024, 1024)) — For most cases,target_size
should be set to the desired height and width of the generated image. If not specified it will default to(height, width)
. Part of SDXL’s micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952. - negative_original_size (
Tuple[int]
, optional, defaults to (1024, 1024)) — To negatively condition the generation process based on a specific image resolution. Part of SDXL’s micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952. For more information, refer to this issue thread: https://github.com/huggingface/diffusers/issues/4208. - negative_crops_coords_top_left (
Tuple[int]
, optional, defaults to (0, 0)) — To negatively condition the generation process based on a specific crop coordinates. Part of SDXL’s micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952. For more information, refer to this issue thread: https://github.com/huggingface/diffusers/issues/4208. - negative_target_size (
Tuple[int]
, optional, defaults to (1024, 1024)) — To negatively condition the generation process based on a target image resolution. It should be as same as thetarget_size
for most cases. Part of SDXL’s micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952. For more information, refer to this issue thread: https://github.com/huggingface/diffusers/issues/4208. - aesthetic_score (
float
, optional, defaults to 6.0) — Used to simulate an aesthetic score of the generated image by influencing the positive text condition. Part of SDXL’s micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952. - negative_aesthetic_score (
float
, optional, defaults to 2.5) — Part of SDXL’s micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952. Can be used to simulate an aesthetic score of the generated image by influencing the negative text condition. - clip_skip (
int
, optional) — Number of layers to be skipped from CLIP while computing the prompt embeddings. A value of 1 means that the output of the pre-final layer will be used for computing the prompt embeddings. - callback_on_step_end (
Callable
, optional) — A function that calls at the end of each denoising steps during the inference. The function is called with the following arguments:callback_on_step_end(self: DiffusionPipeline, step: int, timestep: int, callback_kwargs: Dict)
.callback_kwargs
will include a list of all tensors as specified bycallback_on_step_end_tensor_inputs
. - callback_on_step_end_tensor_inputs (
List
, optional) — The list of tensor inputs for thecallback_on_step_end
function. The tensors specified in the list will be passed ascallback_kwargs
argument. You will only be able to include variables listed in the._callback_tensor_inputs
attribute of your pipeine class.
Returns
StableDiffusionPipelineOutput or tuple
StableDiffusionPipelineOutput if return_dict
is True, otherwise a tuple
containing the output images.
Function invoked when calling the pipeline for generation.
Examples:
>>> # pip install accelerate transformers safetensors diffusers
>>> import torch
>>> import numpy as np
>>> from PIL import Image
>>> from transformers import DPTFeatureExtractor, DPTForDepthEstimation
>>> from diffusers import ControlNetModel, StableDiffusionXLControlNetImg2ImgPipeline, AutoencoderKL
>>> from diffusers.utils import load_image
>>> depth_estimator = DPTForDepthEstimation.from_pretrained("Intel/dpt-hybrid-midas").to("cuda")
>>> feature_extractor = DPTFeatureExtractor.from_pretrained("Intel/dpt-hybrid-midas")
>>> controlnet = ControlNetModel.from_pretrained(
... "diffusers/controlnet-depth-sdxl-1.0-small",
... variant="fp16",
... use_safetensors=True,
... torch_dtype=torch.float16,
... ).to("cuda")
>>> vae = AutoencoderKL.from_pretrained("madebyollin/sdxl-vae-fp16-fix", torch_dtype=torch.float16).to("cuda")
>>> pipe = StableDiffusionXLControlNetImg2ImgPipeline.from_pretrained(
... "stabilityai/stable-diffusion-xl-base-1.0",
... controlnet=controlnet,
... vae=vae,
... variant="fp16",
... use_safetensors=True,
... torch_dtype=torch.float16,
... ).to("cuda")
>>> pipe.enable_model_cpu_offload()
>>> def get_depth_map(image):
... image = feature_extractor(images=image, return_tensors="pt").pixel_values.to("cuda")
... with torch.no_grad(), torch.autocast("cuda"):
... depth_map = depth_estimator(image).predicted_depth
... depth_map = torch.nn.functional.interpolate(
... depth_map.unsqueeze(1),
... size=(1024, 1024),
... mode="bicubic",
... align_corners=False,
... )
... depth_min = torch.amin(depth_map, dim=[1, 2, 3], keepdim=True)
... depth_max = torch.amax(depth_map, dim=[1, 2, 3], keepdim=True)
... depth_map = (depth_map - depth_min) / (depth_max - depth_min)
... image = torch.cat([depth_map] * 3, dim=1)
... image = image.permute(0, 2, 3, 1).cpu().numpy()[0]
... image = Image.fromarray((image * 255.0).clip(0, 255).astype(np.uint8))
... return image
>>> prompt = "A robot, 4k photo"
>>> image = load_image(
... "https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main"
... "/kandinsky/cat.png"
... ).resize((1024, 1024))
>>> controlnet_conditioning_scale = 0.5 # recommended for good generalization
>>> depth_image = get_depth_map(image)
>>> images = pipe(
... prompt,
... image=image,
... control_image=depth_image,
... strength=0.99,
... num_inference_steps=50,
... controlnet_conditioning_scale=controlnet_conditioning_scale,
... ).images
>>> images[0].save(f"robot_cat.png")
Disables the FreeU mechanism if enabled.
Disable sliced VAE decoding. If enable_vae_slicing
was previously enabled, this method will go back to
computing decoding in one step.
Disable tiled VAE decoding. If enable_vae_tiling
was previously enabled, this method will go back to
computing decoding in one step.
enable_freeu
< source >( s1: float s2: float b1: float b2: float )
Parameters
- s1 (
float
) — Scaling factor for stage 1 to attenuate the contributions of the skip features. This is done to mitigate “oversmoothing effect” in the enhanced denoising process. - s2 (
float
) — Scaling factor for stage 2 to attenuate the contributions of the skip features. This is done to mitigate “oversmoothing effect” in the enhanced denoising process. - b1 (
float
) — Scaling factor for stage 1 to amplify the contributions of backbone features. - b2 (
float
) — Scaling factor for stage 2 to amplify the contributions of backbone features.
Enables the FreeU mechanism as in https://arxiv.org/abs/2309.11497.
The suffixes after the scaling factors represent the stages where they are being applied.
Please refer to the official repository for combinations of the values that are known to work well for different pipelines such as Stable Diffusion v1, v2, and Stable Diffusion XL.
Enable sliced VAE decoding. When this option is enabled, the VAE will split the input tensor in slices to compute decoding in several steps. This is useful to save some memory and allow larger batch sizes.
Enable tiled VAE decoding. When this option is enabled, the VAE will split the input tensor into tiles to compute decoding and encoding in several steps. This is useful for saving a large amount of memory and to allow processing larger images.
encode_prompt
< source >( prompt: str prompt_2: Optional = None device: Optional = None num_images_per_prompt: int = 1 do_classifier_free_guidance: bool = True negative_prompt: Optional = None negative_prompt_2: Optional = None prompt_embeds: Optional = None negative_prompt_embeds: Optional = None pooled_prompt_embeds: Optional = None negative_pooled_prompt_embeds: Optional = None lora_scale: Optional = None clip_skip: Optional = None )
Parameters
- prompt (
str
orList[str]
, optional) — prompt to be encoded - prompt_2 (
str
orList[str]
, optional) — The prompt or prompts to be sent to thetokenizer_2
andtext_encoder_2
. If not defined,prompt
is used in both text-encoders device — (torch.device
): torch device - num_images_per_prompt (
int
) — number of images that should be generated per prompt - do_classifier_free_guidance (
bool
) — whether to use classifier free guidance or not - negative_prompt (
str
orList[str]
, optional) — The prompt or prompts not to guide the image generation. If not defined, one has to passnegative_prompt_embeds
instead. Ignored when not using guidance (i.e., ignored ifguidance_scale
is less than1
). - negative_prompt_2 (
str
orList[str]
, optional) — The prompt or prompts not to guide the image generation to be sent totokenizer_2
andtext_encoder_2
. If not defined,negative_prompt
is used in both text-encoders - prompt_embeds (
torch.FloatTensor
, optional) — Pre-generated text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not provided, text embeddings will be generated fromprompt
input argument. - negative_prompt_embeds (
torch.FloatTensor
, optional) — Pre-generated negative text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not provided, negative_prompt_embeds will be generated fromnegative_prompt
input argument. - pooled_prompt_embeds (
torch.FloatTensor
, optional) — Pre-generated pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not provided, pooled text embeddings will be generated fromprompt
input argument. - negative_pooled_prompt_embeds (
torch.FloatTensor
, optional) — Pre-generated negative pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not provided, pooled negative_prompt_embeds will be generated fromnegative_prompt
input argument. - lora_scale (
float
, optional) — A lora scale that will be applied to all LoRA layers of the text encoder if LoRA layers are loaded. - clip_skip (
int
, optional) — Number of layers to be skipped from CLIP while computing the prompt embeddings. A value of 1 means that the output of the pre-final layer will be used for computing the prompt embeddings.
Encodes the prompt into text encoder hidden states.
StableDiffusionXLControlNetInpaintPipeline
class diffusers.StableDiffusionXLControlNetInpaintPipeline
< source >( vae: AutoencoderKL text_encoder: CLIPTextModel text_encoder_2: CLIPTextModelWithProjection tokenizer: CLIPTokenizer tokenizer_2: CLIPTokenizer unet: UNet2DConditionModel controlnet: ControlNetModel scheduler: KarrasDiffusionSchedulers requires_aesthetics_score: bool = False force_zeros_for_empty_prompt: bool = True add_watermarker: Optional = None )
Parameters
- vae (AutoencoderKL) — Variational Auto-Encoder (VAE) Model to encode and decode images to and from latent representations.
- text_encoder (
CLIPTextModel
) — Frozen text-encoder. Stable Diffusion XL uses the text portion of CLIP, specifically the clip-vit-large-patch14 variant. - text_encoder_2 (
CLIPTextModelWithProjection
) — Second frozen text-encoder. Stable Diffusion XL uses the text and pool portion of CLIP, specifically the laion/CLIP-ViT-bigG-14-laion2B-39B-b160k variant. - tokenizer (
CLIPTokenizer
) — Tokenizer of class CLIPTokenizer. - tokenizer_2 (
CLIPTokenizer
) — Second Tokenizer of class CLIPTokenizer. - unet (UNet2DConditionModel) — Conditional U-Net architecture to denoise the encoded image latents.
- scheduler (SchedulerMixin) —
A scheduler to be used in combination with
unet
to denoise the encoded image latents. Can be one of DDIMScheduler, LMSDiscreteScheduler, or PNDMScheduler.
Pipeline for text-to-image generation using Stable Diffusion XL.
This model inherits from DiffusionPipeline. Check the superclass documentation for the generic methods the library implements for all the pipelines (such as downloading or saving, running on a particular device, etc.)
The pipeline also inherits the following loading methods:
- load_lora_weights() for loading LoRA weights
save_lora_weights()
for saving LoRA weights- from_single_file() for loading
.ckpt
files
__call__
< source >( prompt: Union = None prompt_2: Union = None image: Union = None mask_image: Union = None control_image: Union = None height: Optional = None width: Optional = None padding_mask_crop: Optional = None strength: float = 0.9999 num_inference_steps: int = 50 denoising_start: Optional = None denoising_end: Optional = None guidance_scale: float = 5.0 negative_prompt: Union = None negative_prompt_2: Union = None num_images_per_prompt: Optional = 1 eta: float = 0.0 generator: Union = None latents: Optional = None prompt_embeds: Optional = None negative_prompt_embeds: Optional = None pooled_prompt_embeds: Optional = None negative_pooled_prompt_embeds: Optional = None output_type: Optional = 'pil' return_dict: bool = True cross_attention_kwargs: Optional = None controlnet_conditioning_scale: Union = 1.0 guess_mode: bool = False control_guidance_start: Union = 0.0 control_guidance_end: Union = 1.0 guidance_rescale: float = 0.0 original_size: Tuple = None crops_coords_top_left: Tuple = (0, 0) target_size: Tuple = None aesthetic_score: float = 6.0 negative_aesthetic_score: float = 2.5 clip_skip: Optional = None callback_on_step_end: Optional = None callback_on_step_end_tensor_inputs: List = ['latents'] **kwargs ) → ~pipelines.stable_diffusion.StableDiffusionXLPipelineOutput
or tuple
Parameters
- prompt (
str
orList[str]
, optional) — The prompt or prompts to guide the image generation. If not defined, one has to passprompt_embeds
. instead. - prompt_2 (
str
orList[str]
, optional) — The prompt or prompts to be sent to thetokenizer_2
andtext_encoder_2
. If not defined,prompt
is used in both text-encoders - image (
PIL.Image.Image
) —Image
, or tensor representing an image batch which will be inpainted, i.e. parts of the image will be masked out withmask_image
and repainted according toprompt
. - mask_image (
PIL.Image.Image
) —Image
, or tensor representing an image batch, to maskimage
. White pixels in the mask will be repainted, while black pixels will be preserved. Ifmask_image
is a PIL image, it will be converted to a single channel (luminance) before use. If it’s a tensor, it should contain one color channel (L) instead of 3, so the expected shape would be(B, H, W, 1)
. - height (
int
, optional, defaults to self.unet.config.sample_size * self.vae_scale_factor) — The height in pixels of the generated image. - width (
int
, optional, defaults to self.unet.config.sample_size * self.vae_scale_factor) — The width in pixels of the generated image. - padding_mask_crop (
int
, optional, defaults toNone
) — The size of margin in the crop to be applied to the image and masking. IfNone
, no crop is applied to image and mask_image. Ifpadding_mask_crop
is notNone
, it will first find a rectangular region with the same aspect ration of the image and contains all masked area, and then expand that area based onpadding_mask_crop
. The image and mask_image will then be cropped based on the expanded area before resizing to the original image size for inpainting. This is useful when the masked area is small while the image is large and contain information inreleant for inpainging, such as background. - strength (
float
, optional, defaults to 0.9999) — Conceptually, indicates how much to transform the masked portion of the referenceimage
. Must be between 0 and 1.image
will be used as a starting point, adding more noise to it the larger thestrength
. The number of denoising steps depends on the amount of noise initially added. Whenstrength
is 1, added noise will be maximum and the denoising process will run for the full number of iterations specified innum_inference_steps
. A value of 1, therefore, essentially ignores the masked portion of the referenceimage
. Note that in the case ofdenoising_start
being declared as an integer, the value ofstrength
will be ignored. - num_inference_steps (
int
, optional, defaults to 50) — The number of denoising steps. More denoising steps usually lead to a higher quality image at the expense of slower inference. - denoising_start (
float
, optional) — When specified, indicates the fraction (between 0.0 and 1.0) of the total denoising process to be bypassed before it is initiated. Consequently, the initial part of the denoising process is skipped and it is assumed that the passedimage
is a partly denoised image. Note that when this is specified, strength will be ignored. Thedenoising_start
parameter is particularly beneficial when this pipeline is integrated into a “Mixture of Denoisers” multi-pipeline setup, as detailed in Refining the Image Output. - denoising_end (
float
, optional) — When specified, determines the fraction (between 0.0 and 1.0) of the total denoising process to be completed before it is intentionally prematurely terminated. As a result, the returned sample will still retain a substantial amount of noise (ca. final 20% of timesteps still needed) and should be denoised by a successor pipeline that hasdenoising_start
set to 0.8 so that it only denoises the final 20% of the scheduler. The denoising_end parameter should ideally be utilized when this pipeline forms a part of a “Mixture of Denoisers” multi-pipeline setup, as elaborated in Refining the Image Output. - guidance_scale (
float
, optional, defaults to 7.5) — Guidance scale as defined in Classifier-Free Diffusion Guidance.guidance_scale
is defined asw
of equation 2. of Imagen Paper. Guidance scale is enabled by settingguidance_scale > 1
. Higher guidance scale encourages to generate images that are closely linked to the textprompt
, usually at the expense of lower image quality. - negative_prompt (
str
orList[str]
, optional) — The prompt or prompts not to guide the image generation. If not defined, one has to passnegative_prompt_embeds
instead. Ignored when not using guidance (i.e., ignored ifguidance_scale
is less than1
). - negative_prompt_2 (
str
orList[str]
, optional) — The prompt or prompts not to guide the image generation to be sent totokenizer_2
andtext_encoder_2
. If not defined,negative_prompt
is used in both text-encoders - prompt_embeds (
torch.FloatTensor
, optional) — Pre-generated text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not provided, text embeddings will be generated fromprompt
input argument. - negative_prompt_embeds (
torch.FloatTensor
, optional) — Pre-generated negative text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not provided, negative_prompt_embeds will be generated fromnegative_prompt
input argument. - pooled_prompt_embeds (
torch.FloatTensor
, optional) — Pre-generated pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not provided, pooled text embeddings will be generated fromprompt
input argument. - negative_pooled_prompt_embeds (
torch.FloatTensor
, optional) — Pre-generated negative pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not provided, pooled negative_prompt_embeds will be generated fromnegative_prompt
input argument. - num_images_per_prompt (
int
, optional, defaults to 1) — The number of images to generate per prompt. - eta (
float
, optional, defaults to 0.0) — Corresponds to parameter eta (η) in the DDIM paper: https://arxiv.org/abs/2010.02502. Only applies to schedulers.DDIMScheduler, will be ignored for others. - generator (
torch.Generator
, optional) — One or a list of torch generator(s) to make generation deterministic. - latents (
torch.FloatTensor
, optional) — Pre-generated noisy latents, sampled from a Gaussian distribution, to be used as inputs for image generation. Can be used to tweak the same generation with different prompts. If not provided, a latents tensor will ge generated by sampling using the supplied randomgenerator
. - output_type (
str
, optional, defaults to"pil"
) — The output format of the generate image. Choose between PIL:PIL.Image.Image
ornp.array
. - return_dict (
bool
, optional, defaults toTrue
) — Whether or not to return a StableDiffusionPipelineOutput instead of a plain tuple. - cross_attention_kwargs (
dict
, optional) — A kwargs dictionary that if specified is passed along to theAttentionProcessor
as defined underself.processor
in diffusers.models.attention_processor. - original_size (
Tuple[int]
, optional, defaults to (1024, 1024)) — Iforiginal_size
is not the same astarget_size
the image will appear to be down- or upsampled.original_size
defaults to(width, height)
if not specified. Part of SDXL’s micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952. - crops_coords_top_left (
Tuple[int]
, optional, defaults to (0, 0)) —crops_coords_top_left
can be used to generate an image that appears to be “cropped” from the positioncrops_coords_top_left
downwards. Favorable, well-centered images are usually achieved by settingcrops_coords_top_left
to (0, 0). Part of SDXL’s micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952. - target_size (
Tuple[int]
, optional, defaults to (1024, 1024)) — For most cases,target_size
should be set to the desired height and width of the generated image. If not specified it will default to(width, height)
. Part of SDXL’s micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952. - aesthetic_score (
float
, optional, defaults to 6.0) — Used to simulate an aesthetic score of the generated image by influencing the positive text condition. Part of SDXL’s micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952. - negative_aesthetic_score (
float
, optional, defaults to 2.5) — Part of SDXL’s micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952. Can be used to simulate an aesthetic score of the generated image by influencing the negative text condition. - clip_skip (
int
, optional) — Number of layers to be skipped from CLIP while computing the prompt embeddings. A value of 1 means that the output of the pre-final layer will be used for computing the prompt embeddings. - callback_on_step_end (
Callable
, optional) — A function that calls at the end of each denoising steps during the inference. The function is called with the following arguments:callback_on_step_end(self: DiffusionPipeline, step: int, timestep: int, callback_kwargs: Dict)
.callback_kwargs
will include a list of all tensors as specified bycallback_on_step_end_tensor_inputs
. - callback_on_step_end_tensor_inputs (
List
, optional) — The list of tensor inputs for thecallback_on_step_end
function. The tensors specified in the list will be passed ascallback_kwargs
argument. You will only be able to include variables listed in the._callback_tensor_inputs
attribute of your pipeine class.
Returns
~pipelines.stable_diffusion.StableDiffusionXLPipelineOutput
or tuple
~pipelines.stable_diffusion.StableDiffusionXLPipelineOutput
if return_dict
is True, otherwise a
tuple.
tuple. When returning a tuple, the first element is a list with the generated images.
Function invoked when calling the pipeline for generation.
Examples:
>>> # !pip install transformers accelerate
>>> from diffusers import StableDiffusionXLControlNetInpaintPipeline, ControlNetModel, DDIMScheduler
>>> from diffusers.utils import load_image
>>> import numpy as np
>>> import torch
>>> init_image = load_image(
... "https://huggingface.co/datasets/diffusers/test-arrays/resolve/main/stable_diffusion_inpaint/boy.png"
... )
>>> init_image = init_image.resize((1024, 1024))
>>> generator = torch.Generator(device="cpu").manual_seed(1)
>>> mask_image = load_image(
... "https://huggingface.co/datasets/diffusers/test-arrays/resolve/main/stable_diffusion_inpaint/boy_mask.png"
... )
>>> mask_image = mask_image.resize((1024, 1024))
>>> def make_canny_condition(image):
... image = np.array(image)
... image = cv2.Canny(image, 100, 200)
... image = image[:, :, None]
... image = np.concatenate([image, image, image], axis=2)
... image = Image.fromarray(image)
... return image
>>> control_image = make_canny_condition(init_image)
>>> controlnet = ControlNetModel.from_pretrained(
... "diffusers/controlnet-canny-sdxl-1.0", torch_dtype=torch.float16
... )
>>> pipe = StableDiffusionXLControlNetInpaintPipeline.from_pretrained(
... "stabilityai/stable-diffusion-xl-base-1.0", controlnet=controlnet, torch_dtype=torch.float16
... )
>>> pipe.enable_model_cpu_offload()
>>> # generate image
>>> image = pipe(
... "a handsome man with ray-ban sunglasses",
... num_inference_steps=20,
... generator=generator,
... eta=1.0,
... image=init_image,
... mask_image=mask_image,
... control_image=control_image,
... ).images[0]
Disables the FreeU mechanism if enabled.
Disable sliced VAE decoding. If enable_vae_slicing
was previously enabled, this method will go back to
computing decoding in one step.
Disable tiled VAE decoding. If enable_vae_tiling
was previously enabled, this method will go back to
computing decoding in one step.
enable_freeu
< source >( s1: float s2: float b1: float b2: float )
Parameters
- s1 (
float
) — Scaling factor for stage 1 to attenuate the contributions of the skip features. This is done to mitigate “oversmoothing effect” in the enhanced denoising process. - s2 (
float
) — Scaling factor for stage 2 to attenuate the contributions of the skip features. This is done to mitigate “oversmoothing effect” in the enhanced denoising process. - b1 (
float
) — Scaling factor for stage 1 to amplify the contributions of backbone features. - b2 (
float
) — Scaling factor for stage 2 to amplify the contributions of backbone features.
Enables the FreeU mechanism as in https://arxiv.org/abs/2309.11497.
The suffixes after the scaling factors represent the stages where they are being applied.
Please refer to the official repository for combinations of the values that are known to work well for different pipelines such as Stable Diffusion v1, v2, and Stable Diffusion XL.
Enable sliced VAE decoding. When this option is enabled, the VAE will split the input tensor in slices to compute decoding in several steps. This is useful to save some memory and allow larger batch sizes.
Enable tiled VAE decoding. When this option is enabled, the VAE will split the input tensor into tiles to compute decoding and encoding in several steps. This is useful for saving a large amount of memory and to allow processing larger images.
encode_prompt
< source >( prompt: str prompt_2: Optional = None device: Optional = None num_images_per_prompt: int = 1 do_classifier_free_guidance: bool = True negative_prompt: Optional = None negative_prompt_2: Optional = None prompt_embeds: Optional = None negative_prompt_embeds: Optional = None pooled_prompt_embeds: Optional = None negative_pooled_prompt_embeds: Optional = None lora_scale: Optional = None clip_skip: Optional = None )
Parameters
- prompt (
str
orList[str]
, optional) — prompt to be encoded - prompt_2 (
str
orList[str]
, optional) — The prompt or prompts to be sent to thetokenizer_2
andtext_encoder_2
. If not defined,prompt
is used in both text-encoders device — (torch.device
): torch device - num_images_per_prompt (
int
) — number of images that should be generated per prompt - do_classifier_free_guidance (
bool
) — whether to use classifier free guidance or not - negative_prompt (
str
orList[str]
, optional) — The prompt or prompts not to guide the image generation. If not defined, one has to passnegative_prompt_embeds
instead. Ignored when not using guidance (i.e., ignored ifguidance_scale
is less than1
). - negative_prompt_2 (
str
orList[str]
, optional) — The prompt or prompts not to guide the image generation to be sent totokenizer_2
andtext_encoder_2
. If not defined,negative_prompt
is used in both text-encoders - prompt_embeds (
torch.FloatTensor
, optional) — Pre-generated text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not provided, text embeddings will be generated fromprompt
input argument. - negative_prompt_embeds (
torch.FloatTensor
, optional) — Pre-generated negative text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not provided, negative_prompt_embeds will be generated fromnegative_prompt
input argument. - pooled_prompt_embeds (
torch.FloatTensor
, optional) — Pre-generated pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not provided, pooled text embeddings will be generated fromprompt
input argument. - negative_pooled_prompt_embeds (
torch.FloatTensor
, optional) — Pre-generated negative pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not provided, pooled negative_prompt_embeds will be generated fromnegative_prompt
input argument. - lora_scale (
float
, optional) — A lora scale that will be applied to all LoRA layers of the text encoder if LoRA layers are loaded. - clip_skip (
int
, optional) — Number of layers to be skipped from CLIP while computing the prompt embeddings. A value of 1 means that the output of the pre-final layer will be used for computing the prompt embeddings.
Encodes the prompt into text encoder hidden states.
StableDiffusionPipelineOutput
class diffusers.pipelines.stable_diffusion.StableDiffusionPipelineOutput
< source >( images: Union nsfw_content_detected: Optional )
Parameters
- images (
List[PIL.Image.Image]
ornp.ndarray
) — List of denoised PIL images of lengthbatch_size
or NumPy array of shape(batch_size, height, width, num_channels)
. - nsfw_content_detected (
List[bool]
) — List indicating whether the corresponding generated image contains “not-safe-for-work” (nsfw) content orNone
if safety checking could not be performed.
Output class for Stable Diffusion pipelines.