TSAI_S20 / app.py
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import gradio as gr
from base64 import b64encode
import numpy
import torch
from diffusers import AutoencoderKL, LMSDiscreteScheduler, UNet2DConditionModel, StableDiffusionPipeline
# For video display:
from matplotlib import pyplot as plt
from pathlib import Path
from PIL import Image
from torch import autocast
from torchvision import transforms as tfms
from tqdm.auto import tqdm
from transformers import CLIPTextModel, CLIPTokenizer, logging
import os
torch.manual_seed(1)
# Supress some unnecessary warnings when loading the CLIPTextModel
logging.set_verbosity_error()
# Set device
torch_device = "cuda" if torch.cuda.is_available() else "mps" if torch.backends.mps.is_available() else "cpu"
if "mps" == torch_device: os.environ['PYTORCH_ENABLE_MPS_FALLBACK'] = "1"
model_nm = "CompVis/stable-diffusion-v1-4"
output_dir="sd-concept-output"
pipe = StableDiffusionPipeline.from_pretrained(model_nm).to(torch_device)
# Load the autoencoder model which will be used to decode the latents into image space.
vae = pipe.vae
tokenizer = pipe.tokenizer
# Load the tokenizer and text encoder to tokenize and encode the text.
#tokenizer = CLIPTokenizer.from_pretrained("openai/clip-vit-large-patch14", torch_dtype=torch.float16)
text_encoder =pipe.text_encoder
# The UNet model for generating the latents.
unet = pipe.unet
# The noise scheduler
scheduler = LMSDiscreteScheduler(beta_start=0.00085, beta_end=0.012, beta_schedule="scaled_linear", num_train_timesteps=1000)
# To the GPU we go!
vae = vae.to(torch_device)
text_encoder = text_encoder.to(torch_device)
unet = unet.to(torch_device);
pipe.load_textual_inversion("sd-concepts-library/madhubani-art")
pipe.load_textual_inversion("sd-concepts-library/line-art")
pipe.load_textual_inversion("sd-concepts-library/cat-toy")
pipe.load_textual_inversion("sd-concepts-library/concept-art")
def pil_to_latent(input_im):
# Single image -> single latent in a batch (so size 1, 4, 64, 64)
with torch.no_grad():
latent = vae.encode(tfms.ToTensor()(input_im).unsqueeze(0).to(torch_device)*2-1) # Note scaling
return 0.18215 * latent.latent_dist.sample()
def latents_to_pil(latents):
# bath of latents -> list of images
latents = (1 / 0.18215) * latents
with torch.no_grad():
image = vae.decode(latents).sample
image = (image / 2 + 0.5).clamp(0, 1)
image = image.detach().cpu().permute(0, 2, 3, 1).numpy()
images = (image * 255).round().astype("uint8")
pil_images = [Image.fromarray(image) for image in images]
return pil_images
# Prep Scheduler
def set_timesteps(scheduler, num_inference_steps):
scheduler.set_timesteps(num_inference_steps)
scheduler.timesteps = scheduler.timesteps.to(torch.float32) # minor fix to ensure MPS compatibility, fixed in diffusers
def saturation_loss(images):
# Calculate saturation for each pixel in the input image tensor
max_vals, _ = torch.max(images, dim=1, keepdim=True)
min_vals, _ = torch.min(images, dim=1, keepdim=True)
saturation = (max_vals - min_vals) / max_vals.clamp(min=1e-7) # Avoid division by zero
# Calculate mean saturation across the image
mean_saturation = torch.mean(saturation, dim=(2, 3)) # Average over width and height
# Calculate the loss as the negative mean saturation (proportional to saturation)
#loss = torch.abs(saturation - 0.9).mean()
return mean_saturation/10000
def generateImage(prompt, lossScale):
#prompt = 'a puppy in <cat-toy> style' #@param
height = 512 # default height of Stable Diffusion
width = 512 # default width of Stable Diffusion
num_inference_steps = 200 #@param # Number of denoising steps
guidance_scale = 8 #@param # Scale for classifier-free guidance
generator = torch.manual_seed(32) # Seed generator to create the inital latent noise
batch_size = 1
saturation_loss_Scale = lossScale #@param
# Prep text
text_input = tokenizer([prompt], padding="max_length", max_length=tokenizer.model_max_length, truncation=True, return_tensors="pt")
with torch.no_grad():
text_embeddings = text_encoder(text_input.input_ids.to(torch_device))[0]
# And the uncond. input as before:
max_length = text_input.input_ids.shape[-1]
uncond_input = tokenizer(
[""] * batch_size, padding="max_length", max_length=max_length, return_tensors="pt"
)
with torch.no_grad():
uncond_embeddings = text_encoder(uncond_input.input_ids.to(torch_device))[0]
text_embeddings = torch.cat([uncond_embeddings, text_embeddings])
# Prep Scheduler
set_timesteps(scheduler, num_inference_steps)
# Prep latents
latents = torch.randn(
(batch_size, unet.in_channels, height // 8, width // 8),
generator=generator,
)
latents = latents.to(torch_device)
latents = latents * scheduler.init_noise_sigma
# Loop
for i, t in tqdm(enumerate(scheduler.timesteps), total=len(scheduler.timesteps)):
# expand the latents if we are doing classifier-free guidance to avoid doing two forward passes.
latent_model_input = torch.cat([latents] * 2)
sigma = scheduler.sigmas[i]
latent_model_input = scheduler.scale_model_input(latent_model_input, t)
# predict the noise residual
with torch.no_grad():
noise_pred = unet(latent_model_input, t, encoder_hidden_states=text_embeddings)["sample"]
# perform CFG
noise_pred_uncond, noise_pred_text = noise_pred.chunk(2)
noise_pred = noise_pred_uncond + guidance_scale * (noise_pred_text - noise_pred_uncond)
#### ADDITIONAL GUIDANCE ###
if i%5 == 0:
# Requires grad on the latents
latents = latents.detach().requires_grad_()
# Get the predicted x0:
# latents_x0 = latents - sigma * noise_pred
latents_x0 = scheduler.step(noise_pred, t, latents).pred_original_sample
scheduler._step_index = scheduler._step_index - 1
# Decode to image space
denoised_images = vae.decode((1 / 0.18215) * latents_x0).sample / 2 + 0.5 # range (0, 1)
# Calculate loss
loss = saturation_loss(denoised_images) * saturation_loss_Scale
#loss = loss.detach().requires_grad_()
#print('loss.grad_fn = {}'.format(grad_fn))
# Occasionally print it out
if i%10==0:
print(i, 'loss:', loss.item())
# Get gradient
cond_grad = torch.autograd.grad(loss, latents)[0]
# Modify the latents based on this gradient
latents = latents.detach() - cond_grad * sigma**2
# Now step with scheduler
latents = scheduler.step(noise_pred, t, latents).prev_sample
custom_loss_image = latents_to_pil(latents)[0]
return custom_loss_image
def inference(imgText, style, customLoss="no"):
prompt = f'a {imgText} in <{style}> style'
if (customLoss == "yes") :
outImage = generateImage(prompt, 2)
return outImage
else:
outImage = generateImage(prompt, 0)
return outImage
title = "TSAI S20 Assignment: Use a pretrained Sstable Diffusion model and give a demo on its workig"
description = "A simple Gradio interface that accepts a text and style, and generated an image using stable diffusion pipeline"
examples = [["puppy","cat-toy","yes"]]
demo = gr.Interface(
inference,
inputs = [gr.Textbox("Enter an image you want to generate"),
gr.Dropdown(["madhubani-art", "line-art", "cat-toy","concept-art"], label="Choose your style"),
gr.Radio(["yes", "no"], label="Add custom saturation loss?")
],
outputs = [gr.Image(shape=(512, 512), label="Generated Image")],
title = title,
description = description,
examples = examples,
)
demo.launch()