|
--- |
|
license: cc-by-nc-4.0 |
|
library_name: diffusers |
|
base_model: stabilityai/stable-diffusion-3-medium |
|
tags: |
|
- lora |
|
- text-to-image |
|
inference: False |
|
--- |
|
# ⚡ Flash Diffusion: FlashSD3 ⚡ |
|
|
|
|
|
Flash Diffusion is a diffusion distillation method proposed in [Flash Diffusion: Accelerating Any Conditional |
|
Diffusion Model for Few Steps Image Generation](http://arxiv.org/abs/2406.02347) *by Clément Chadebec, Onur Tasar, Eyal Benaroche, and Benjamin Aubin.* |
|
This model is a **90.4M** LoRA distilled version of [SD3](https://huggingface.co/stabilityai/stable-diffusion-3-medium) model that is able to generate 1024x1024 images in **4 steps**. |
|
See our [live demo](https://huggingface.co/spaces/jasperai/flash-sd3) and official [Github repo](https://github.com/gojasper/flash-diffusion). |
|
|
|
|
|
<p align="center"> |
|
<img style="width:700px;" src="assets/flash_sd3.png"> |
|
</p> |
|
|
|
# How to use? |
|
|
|
The model can be used using the `StableDiffusion3Pipeline` from `diffusers` library directly. It can allow reducing the number of required sampling steps to **4 steps**. |
|
|
|
First, you need to install a specific version of `diffusers` by runniung |
|
|
|
```bash |
|
pip install git+https://github.com/initml/diffusers.git@clement/feature/flash_sd3 |
|
``` |
|
|
|
Then, you can ru the following to generate an image |
|
|
|
```python |
|
import torch |
|
from diffusers import StableDiffusion3Pipeline, SD3Transformer2DModel, FlashFlowMatchEulerDiscreteScheduler |
|
from peft import PeftModel |
|
|
|
# Load LoRA |
|
transformer = SD3Transformer2DModel.from_pretrained( |
|
"stabilityai/stable-diffusion-3-medium-diffusers", |
|
subfolder="transformer", |
|
torch_dtype=torch.float16, |
|
) |
|
transformer = PeftModel.from_pretrained(transformer, "jasperai/flash-sd3") |
|
|
|
|
|
# Pipeline |
|
pipe = StableDiffusion3Pipeline.from_pretrained( |
|
"stabilityai/stable-diffusion-3-medium-diffusers", |
|
transformer=transformer, |
|
torch_dtype=torch.float16, |
|
text_encoder_3=None, |
|
tokenizer_3=None |
|
) |
|
|
|
# Scheduler |
|
pipe.scheduler = FlashFlowMatchEulerDiscreteScheduler.from_pretrained( |
|
"stabilityai/stable-diffusion-3-medium-diffusers", |
|
subfolder="scheduler", |
|
) |
|
|
|
pipe.to("cuda") |
|
|
|
prompt = "A raccoon trapped inside a glass jar full of colorful candies, the background is steamy with vivid colors." |
|
|
|
image = pipe(prompt, num_inference_steps=4, guidance_scale=0).images[0] |
|
``` |
|
<p align="center"> |
|
<img style="width:400px;" src="assets/raccoon.png"> |
|
</p> |
|
|
|
## Citation |
|
If you find this work useful or use it in your research, please consider citing us |
|
|
|
```bibtex |
|
@misc{chadebec2024flash, |
|
title={Flash Diffusion: Accelerating Any Conditional Diffusion Model for Few Steps Image Generation}, |
|
author={Clement Chadebec and Onur Tasar and Eyal Benaroche and Benjamin Aubin}, |
|
year={2024}, |
|
eprint={2406.02347}, |
|
archivePrefix={arXiv}, |
|
primaryClass={cs.CV} |
|
} |
|
``` |
|
|
|
## License |
|
This model is released under the the Creative Commons BY-NC license. |