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--- |
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language: en |
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pipeline_tag: text-to-image |
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tags: |
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- stable-diffusion |
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- stable-diffusion-xl |
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- text-to-image |
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- diffusers |
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- lora |
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inference: true |
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base_model: |
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- stabilityai/stable-diffusion-3.5-large |
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--- |
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# SimpleTuner-stable-diffusion-3.5-large-ZukyTheDog |
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This is a custom-trained LoRA model for Stable Diffusion 3.5, focused on generating images of ZukyTheDog. |
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## Model Description |
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This model is a LoRA fine-tune of Stable Diffusion 3.5, trained to generate high-quality images of ZukyTheDog. It can be used to create detailed and accurate representations in various styles and settings. |
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## Base Model |
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- Base Model: Stable Diffusion 3.5 |
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## How to Use |
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You can use this model in two ways: |
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### 1. Using Hugging Face Inference API |
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Simply input your prompt in the inference API tab above and click "Compute". |
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Example prompts: |
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- "A photo of ZukyTheDog sitting in a garden" |
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- "ZukyTheDog playing with a ball, high quality, detailed" |
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- "Portrait of ZukyTheDog in a natural setting" |
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### 2. Using Python Code |
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```python |
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from diffusers import StableDiffusion3Pipeline |
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import torch |
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pipe = StableDiffusion3Pipeline.from_pretrained("stabilityai/stable-diffusion-3.5-large", torch_dtype=torch.bfloat16) |
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pipe.to("cuda") |
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pipe.load_lycoris_weights("CarlosRiverMe/SimpleTuner-stable-diffusion-3.5-large-ZukyTheDog", weight_name="pytorch_lora_weights_final.safetensors") |
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prompt = "A photo of ZukyTheDog sitting in a garden" |
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image = pipe(prompt).images[0] |
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image.save("generated_image.png") |
