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import torch | |
from transformers import CLIPTextModel, CLIPTokenizer | |
from diffusers import AutoencoderKL, UNet2DConditionModel, LMSDiscreteScheduler | |
from tqdm.auto import tqdm | |
from torch import autocast | |
from PIL import Image | |
from matplotlib import pyplot as plt | |
import numpy | |
from torchvision import transforms as tfms | |
import shutil | |
# For video display: | |
import cv2 | |
from IPython.display import HTML | |
from base64 import b64encode | |
import os | |
from utils import color_loss,latents_to_pil,pil_to_latent,sketch_loss | |
# Set device | |
torch_device = "cpu" | |
vae = AutoencoderKL.from_pretrained("CompVis/stable-diffusion-v1-4", subfolder="vae") | |
# Load the tokenizer and text encoder to tokenize and encode the text. | |
tokenizer = CLIPTokenizer.from_pretrained("openai/clip-vit-large-patch14") | |
text_encoder = CLIPTextModel.from_pretrained("openai/clip-vit-large-patch14") | |
# The UNet model for generating the latents. | |
unet = UNet2DConditionModel.from_pretrained("CompVis/stable-diffusion-v1-4", subfolder="unet") | |
# The noise scheduler | |
scheduler = LMSDiscreteScheduler(beta_start=0.00085, beta_end=0.012, beta_schedule="scaled_linear", num_train_timesteps=1000) | |
vae = vae.to(torch_device) | |
text_encoder = text_encoder.to(torch_device) | |
unet = unet.to(torch_device) | |
scheduler.set_timesteps(15) | |
def generate_mixed_image(prompt1, prompt2,noised_image=False): | |
mix_factor = 0.4 #@param | |
height = 512 # default height of Stable Diffusion | |
width = 512 # default width of Stable Diffusion | |
num_inference_steps = 50 #@param # Number of denoising steps | |
guidance_scale = 8 # Scale for classifier-free guidance | |
generator = torch.manual_seed(32) # Seed generator to create the inital latent noise | |
batch_size = 1 | |
# Prep text | |
# Embed both prompts | |
text_input1 = tokenizer([prompt1], padding="max_length", max_length=tokenizer.model_max_length, truncation=True, return_tensors="pt") | |
with torch.no_grad(): | |
text_embeddings1 = text_encoder(text_input1.input_ids.to(torch_device))[0] | |
text_input2 = tokenizer([prompt2], padding="max_length", max_length=tokenizer.model_max_length, truncation=True, return_tensors="pt") | |
with torch.no_grad(): | |
text_embeddings2 = text_encoder(text_input2.input_ids.to(torch_device))[0] | |
# Take the average | |
text_embeddings = (text_embeddings1*mix_factor + \ | |
text_embeddings2*(1-mix_factor)) | |
# And the uncond. input as before: | |
max_length = max(text_input1.input_ids.shape[-1],text_input2.input_ids.shape[-1]) | |
uncond_input = tokenizer( | |
[""] * batch_size, padding="max_length", max_length=max_length, return_tensors="pt" | |
) | |
with torch.no_grad(): | |
uncond_embeddings = text_encoder(uncond_input.input_ids.to(torch_device))[0] | |
text_embeddings = torch.cat([uncond_embeddings, text_embeddings]) | |
# Prep Scheduler | |
scheduler.set_timesteps(num_inference_steps) | |
# Prep latents | |
latents = torch.randn( | |
(batch_size, unet.in_channels, height // 8, width // 8), | |
generator=generator, | |
) | |
latents = latents.to(torch_device) | |
latents = latents * scheduler.sigmas[0] # Need to scale to match k | |
# Loop | |
with autocast("cuda"): | |
for i, t in tqdm(enumerate(scheduler.timesteps)): | |
# expand the latents if we are doing classifier-free guidance to avoid doing two forward passes. | |
latent_model_input = torch.cat([latents] * 2) | |
sigma = scheduler.sigmas[i] | |
latent_model_input = latent_model_input / ((sigma**2 + 1) ** 0.5) | |
# predict the noise residual | |
with torch.no_grad(): | |
noise_pred = unet(latent_model_input, t, encoder_hidden_states=text_embeddings)["sample"] | |
# perform guidance | |
noise_pred_uncond, noise_pred_text = noise_pred.chunk(2) | |
noise_pred = noise_pred_uncond + guidance_scale * (noise_pred_text - noise_pred_uncond) | |
# compute the previous noisy sample x_t -> x_t-1 | |
latents = scheduler.step(noise_pred, i, latents)["prev_sample"] | |
if noised_image: | |
output = generate_noised_version_of_image(latents_to_pil(latents,vae)[0]) | |
else: | |
output = latents_to_pil(latents,vae)[0] | |
return output | |
def generate_image(prompt,color_postprocessing=False,postporcessing_color=None,color_loss_scale=40,noised_image=False): | |
#@title Store the predicted outputs and next frame for later viewing | |
#prompt = 'A campfire (oil on canvas)' # | |
height = 512 # default height of Stable Diffusion | |
width = 512 # default width of Stable Diffusion | |
num_inference_steps = 50 # # Number of denoising steps | |
guidance_scale = 8 # # Scale for classifier-free guidance | |
generator = torch.manual_seed(32) # Seed generator to create the inital latent noise | |
batch_size = 1 | |
# Define the directory name | |
directory_name = "steps" | |
# Check if the directory exists, and if so, delete it | |
if os.path.exists(directory_name): | |
shutil.rmtree(directory_name) | |
#Create the directory | |
os.makedirs(directory_name) | |
# Prep text | |
text_input = tokenizer([prompt], padding="max_length", max_length=tokenizer.model_max_length, truncation=True, return_tensors="pt") | |
with torch.no_grad(): | |
text_embeddings = text_encoder(text_input.input_ids.to(torch_device))[0] | |
# And the uncond. input as before: | |
max_length = text_input.input_ids.shape[-1] | |
uncond_input = tokenizer( | |
[""] * batch_size, padding="max_length", max_length=max_length, return_tensors="pt" | |
) | |
with torch.no_grad(): | |
uncond_embeddings = text_encoder(uncond_input.input_ids.to(torch_device))[0] | |
text_embeddings = torch.cat([uncond_embeddings, text_embeddings]) | |
# Prep Scheduler | |
scheduler.set_timesteps(num_inference_steps) | |
# Prep latents | |
latents = torch.randn( | |
(batch_size, unet.in_channels, height // 8, width // 8), | |
generator=generator, | |
) | |
latents = latents.to(torch_device) | |
latents = latents * scheduler.sigmas[0] # Need to scale to match k | |
# Loop | |
with autocast("cuda"): | |
for i, t in tqdm(enumerate(scheduler.timesteps)): | |
# expand the latents if we are doing classifier-free guidance to avoid doing two forward passes. | |
latent_model_input = torch.cat([latents] * 2) | |
sigma = scheduler.sigmas[i] | |
latent_model_input = latent_model_input / ((sigma**2 + 1) ** 0.5) | |
# predict the noise residual | |
with torch.no_grad(): | |
noise_pred = unet(latent_model_input, t, encoder_hidden_states=text_embeddings)["sample"] | |
# perform CFG | |
noise_pred_uncond, noise_pred_text = noise_pred.chunk(2) | |
noise_pred = noise_pred_uncond + guidance_scale * (noise_pred_text - noise_pred_uncond) | |
#### ADDITIONAL GUIDANCE ### | |
# Requires grad on the latents | |
if color_postprocessing: | |
latents = latents.detach().requires_grad_() | |
# Get the predicted x0: | |
latents_x0 = latents - sigma * noise_pred | |
# Decode to image space | |
denoised_images = vae.decode((1 / 0.18215) * latents_x0) / 2 + 0.5 # (0, 1) | |
# Calculate loss | |
#loss = sketch_loss(denoised_images) * color_loss_scale | |
loss = color_loss(denoised_images,postporcessing_color) * color_loss_scale | |
if i%10==0: | |
print(i, 'loss:', loss.item()) | |
# Get gradient | |
cond_grad = -torch.autograd.grad(loss, latents)[0] | |
# Modify the latents based on this gradient | |
latents = latents.detach() + cond_grad * sigma**2 | |
### And saving as before ### | |
# Get the predicted x0: | |
latents_x0 = latents - sigma * noise_pred | |
im_t0 = latents_to_pil(latents_x0,vae)[0] | |
# And the previous noisy sample x_t -> x_t-1 | |
latents = scheduler.step(noise_pred, i, latents)["prev_sample"] | |
im_next = latents_to_pil(latents,vae)[0] | |
# Combine the two images and save for later viewing | |
im = Image.new('RGB', (1024, 512)) | |
im.paste(im_next, (0, 0)) | |
im.paste(im_t0, (512, 0)) | |
im.save(f'steps/{i:04}.jpeg') | |
else: | |
latents = scheduler.step(noise_pred, i, latents)["prev_sample"] | |
if noised_image: | |
output = generate_noised_version_of_image(latents_to_pil(latents,vae)[0]) | |
else: | |
output = latents_to_pil(latents,vae)[0] | |
return output | |
def progress_video(prompt): | |
pil_image = generate_image(prompt) | |
# Generate a list of image file paths (replace with your own logic) | |
num_frames = len(os.listdir("steps/")) | |
image_files = [f"steps/{i:04d}.jpeg" for i in range(1, num_frames + 1)] | |
# Read the first image to get its size (assuming all images have the same size) | |
first_image = cv2.imread({image_files[0]}) | |
height, width, _ = first_image.shape | |
# Define the output video writer | |
fourcc = cv2.VideoWriter_fourcc(*'mp4v') # Codec for MP4 | |
out = cv2.VideoWriter('out.mp4', fourcc, 12, (width, height)) | |
for image_file in image_files: | |
frame = cv2.imread(image_file) | |
out.write(frame) | |
out.release() | |
return "out.mp4" | |
def generate_noised_version_of_image(pil_image): | |
# View a noised version | |
encoded = pil_to_latent(pil_image,vae) | |
noise = torch.randn_like(encoded) # Random noise | |
timestep = 150 # i.e. equivalent to that at 150/1000 training steps | |
encoded_and_noised = scheduler.add_noise(encoded, noise, timestep) | |
return latents_to_pil(encoded_and_noised,vae)[0] # Display | |
# if __name__ == "__main__": | |
# prompt = 'A campfire (oil on canvas)' | |
# color_loss_scale = 40 | |
# color_postprocessing = False | |
# pil_image = generate_mixed_image("a dog", "a cat") | |
# #pil_image = generate_image(prompt,color_postprocessing,color_loss_scale) | |
# #pil_image = generate_noised_version_of_image(Image.open('output.png').resize((512, 512))) | |
# pil_image.save("output1.png") | |
if __name__ == "__main__": | |
progress_video("lol") |