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| from io import BytesIO | |
| import torch | |
| import numpy as np | |
| from PIL import Image | |
| from einops import rearrange | |
| from torch import autocast | |
| from contextlib import nullcontext | |
| import requests | |
| import functools | |
| from ldm.models.diffusion.ddim import DDIMSampler | |
| from ldm.models.diffusion.plms import PLMSSampler | |
| from ldm.extras import load_model_from_config, load_training_dir | |
| import clip | |
| from PIL import Image | |
| from huggingface_hub import hf_hub_download | |
| ckpt = hf_hub_download(repo_id="lambdalabs/image-mixer", filename="image-mixer-pruned.ckpt") | |
| config = hf_hub_download(repo_id="lambdalabs/image-mixer", filename="image-mixer-config.yaml") | |
| device = "cuda:0" | |
| model = load_model_from_config(config, ckpt, device=device, verbose=False) | |
| model = model.to(device).half() | |
| clip_model, preprocess = clip.load("ViT-L/14", device=device) | |
| n_inputs = 5 | |
| torch.cuda.empty_cache() | |
| def get_url_im(t): | |
| user_agent = {'User-agent': 'gradio-app'} | |
| response = requests.get(t, headers=user_agent) | |
| return Image.open(BytesIO(response.content)) | |
| def get_im_c(im_path, clip_model): | |
| # im = Image.open(im_path).convert("RGB") | |
| prompts = preprocess(im_path).to(device).unsqueeze(0) | |
| return clip_model.encode_image(prompts).float() | |
| def get_txt_c(txt, clip_model): | |
| text = clip.tokenize([txt,]).to(device) | |
| return clip_model.encode_text(text) | |
| def get_txt_diff(txt1, txt2, clip_model): | |
| return get_txt_c(txt1, clip_model) - get_txt_c(txt2, clip_model) | |
| def to_im_list(x_samples_ddim): | |
| x_samples_ddim = torch.clamp((x_samples_ddim + 1.0) / 2.0, min=0.0, max=1.0) | |
| ims = [] | |
| for x_sample in x_samples_ddim: | |
| x_sample = 255. * rearrange(x_sample.cpu().numpy(), 'c h w -> h w c') | |
| ims.append(Image.fromarray(x_sample.astype(np.uint8))) | |
| return ims | |
| def sample(sampler, model, c, uc, scale, start_code, h=512, w=512, precision="autocast",ddim_steps=50): | |
| ddim_eta=0.0 | |
| precision_scope = autocast if precision=="autocast" else nullcontext | |
| with precision_scope("cuda"): | |
| shape = [4, h // 8, w // 8] | |
| samples_ddim, _ = sampler.sample(S=ddim_steps, | |
| conditioning=c, | |
| batch_size=c.shape[0], | |
| shape=shape, | |
| verbose=False, | |
| unconditional_guidance_scale=scale, | |
| unconditional_conditioning=uc, | |
| eta=ddim_eta, | |
| x_T=start_code) | |
| x_samples_ddim = model.decode_first_stage(samples_ddim) | |
| return to_im_list(x_samples_ddim) | |
| def run(*args): | |
| inps = [] | |
| for i in range(0, len(args)-4, n_inputs): | |
| inps.append(args[i:i+n_inputs]) | |
| scale, n_samples, seed, steps = args[-4:] | |
| h = w = 640 | |
| sampler = DDIMSampler(model) | |
| # sampler = PLMSSampler(model) | |
| torch.manual_seed(seed) | |
| start_code = torch.randn(n_samples, 4, h//8, w//8, device=device) | |
| conds = [] | |
| for b, t, im, s in zip(*inps): | |
| if b == "Image": | |
| this_cond = s*get_im_c(im, clip_model) | |
| elif b == "Text/URL": | |
| if t.startswith("http"): | |
| im = get_url_im(t) | |
| this_cond = s*get_im_c(im, clip_model) | |
| else: | |
| this_cond = s*get_txt_c(t, clip_model) | |
| else: | |
| this_cond = torch.zeros((1, 768), device=device) | |
| conds.append(this_cond) | |
| conds = torch.cat(conds, dim=0).unsqueeze(0) | |
| conds = conds.tile(n_samples, 1, 1) | |
| ims = sample(sampler, model, conds, 0*conds, scale, start_code, ddim_steps=steps) | |
| # return make_row(ims) | |
| # Clear GPU memory cache so less likely to OOM | |
| torch.cuda.empty_cache() | |
| return ims | |
| import gradio as gr | |
| from functools import partial | |
| from itertools import chain | |
| def change_visible(txt1, im1, val): | |
| outputs = {} | |
| if val == "Image": | |
| outputs[im1] = gr.update(visible=True) | |
| outputs[txt1] = gr.update(visible=False) | |
| elif val == "Text/URL": | |
| outputs[im1] = gr.update(visible=False) | |
| outputs[txt1] = gr.update(visible=True) | |
| elif val == "Nothing": | |
| outputs[im1] = gr.update(visible=False) | |
| outputs[txt1] = gr.update(visible=False) | |
| return outputs | |
| with gr.Blocks(title="Image Mixer", css=".gr-box {border-color: #8136e2}") as demo: | |
| gr.Markdown("") | |
| gr.Markdown( | |
| """ | |
| # Image Mixer | |
| _Created by [Justin Pinkney](https://www.justinpinkney.com) at [Lambda Labs](https://lambdalabs.com/)_ | |
| To skip the queue you can try it on <a href="https://cloud.lambdalabs.com/demos/lambda-ml/image-mixer-demo" style="display:inline-block;position: relative;"><img style="margin-top: 0;margin-bottom: 0;margin-left: .25em;" src="https://img.shields.io/badge/-Lambda%20Cloud-blueviolet"></a> | |
| ### __Provide one or more images to be mixed together by a fine-tuned Stable Diffusion model (see tips and advice below👇).__ | |
|  | |
| """) | |
| btns = [] | |
| txts = [] | |
| ims = [] | |
| strengths = [] | |
| with gr.Row(): | |
| for i in range(n_inputs): | |
| with gr.Box(): | |
| with gr.Column(): | |
| btn1 = gr.Radio( | |
| choices=["Image", "Text/URL", "Nothing"], | |
| label=f"Input {i} type", | |
| interactive=True, | |
| value="Nothing", | |
| ) | |
| txt1 = gr.Textbox(label="Text or Image URL", visible=False, interactive=True) | |
| im1 = gr.Image(label="Image", interactive=True, visible=False, type="pil") | |
| strength = gr.Slider(label="Strength", minimum=0, maximum=5, step=0.05, value=1, interactive=True) | |
| fn = partial(change_visible, txt1, im1) | |
| btn1.change(fn=fn, inputs=[btn1], outputs=[txt1, im1], queue=False) | |
| btns.append(btn1) | |
| txts.append(txt1) | |
| ims.append(im1) | |
| strengths.append(strength) | |
| with gr.Row(): | |
| cfg_scale = gr.Slider(label="CFG scale", value=3, minimum=1, maximum=10, step=0.5) | |
| n_samples = gr.Slider(label="Num samples", value=1, minimum=1, maximum=1, step=1) | |
| seed = gr.Slider(label="Seed", value=0, minimum=0, maximum=10000, step=1) | |
| steps = gr.Slider(label="Steps", value=30, minimum=10, maximum=100, step=5) | |
| with gr.Row(): | |
| submit = gr.Button("Generate") | |
| output = gr.Gallery().style(grid=[1,2], height="640px") | |
| inps = list(chain(btns, txts, ims, strengths)) | |
| inps.extend([cfg_scale,n_samples,seed, steps,]) | |
| submit.click(fn=run, inputs=inps, outputs=[output]) | |
| ex = gr.Examples([ | |
| [ | |
| "Image", "Image", "Text/URL", "Nothing", "Nothing", | |
| "","","central symmetric figure detailed artwork","","", | |
| "gainsborough.jpeg","blonder.jpeg","blonder.jpeg","blonder.jpeg","blonder.jpeg", | |
| 1,1.35,1.4,1,1, | |
| 3.0, 1, 0, 30, | |
| ], | |
| [ | |
| "Image", "Image", "Text/URL", "Image", "Nothing", | |
| "","","flowers","","", | |
| "ex2-1.jpeg","ex2-2.jpeg","blonder.jpeg","ex2-3.jpeg","blonder.jpeg", | |
| 1,1,1.5,1.25,1, | |
| 3.0, 1, 0, 30, | |
| ], | |
| [ | |
| "Image", "Image", "Image", "Nothing", "Nothing", | |
| "","","","","", | |
| "ex1-1.jpeg","ex1-2.jpeg","ex1-3.jpeg","blonder.jpeg","blonder.jpeg", | |
| 1.1,1,1.4,1,1, | |
| 3.0, 1, 0, 30, | |
| ], | |
| ], | |
| fn=run, inputs=inps, outputs=[output], cache_examples=True) | |
| gr.Markdown( | |
| """ | |
| ## Tips | |
| - You can provide between 1 and 5 inputs, these can either be an uploaded image a text prompt or a url to an image file. | |
| - The order of the inputs shouldn't matter, any images will be centre cropped before use. | |
| - Each input has an individual strength parameter which controls how big an influence it has on the output. | |
| - The model was not trained using text and can not interpret complex text prompts. | |
| - Using only text prompts doesn't work well, make sure there is at least one image or URL to an image. | |
| - The parameters on the bottom row such as cfg scale do the same as for a normal Stable Diffusion model. | |
| - Balancing the different inputs requires tweaking of the strengths, I suggest getting the right balance for a small number of samples and with few steps until you're | |
| happy with the result then increase the steps for better quality. | |
| - Outputs are 640x640 by default. | |
| - If you want to run locally see the instruction on the [Model Card](https://huggingface.co/lambdalabs/image-mixer). | |
| ## How does this work? | |
| This model is based on the [Stable Diffusion Image Variations model](https://huggingface.co/lambdalabs/sd-image-variations-diffusers) | |
| but it has been fined tuned to take multiple CLIP image embeddings. During training, up to 5 random crops were taken from the training images and | |
| the CLIP image embeddings were computed, these were then concatenated and used as the conditioning for the model. At inference time we can combine the image | |
| embeddings from multiple images to mix their concepts (and we can also use the text encoder to add text concepts too). | |
| The model was trained on a subset of LAION Improved Aesthetics at a resolution of 640x640 and was trained using 8xA100 GPUs on [Lambda GPU Cloud](https://lambdalabs.com/service/gpu-cloud). | |
| """) | |
| demo.launch() | |