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log_dir: "Models/LibriTTS" | |
first_stage_path: "first_stage.pth" | |
save_freq: 1 | |
log_interval: 10 | |
device: "cuda" | |
epochs_1st: 50 # number of epochs for first stage training (pre-training) | |
epochs_2nd: 30 # number of peochs for second stage training (joint training) | |
batch_size: 16 | |
max_len: 300 # maximum number of frames | |
pretrained_model: "" | |
second_stage_load_pretrained: true # set to true if the pre-trained model is for 2nd stage | |
load_only_params: false # set to true if do not want to load epoch numbers and optimizer parameters | |
F0_path: "Utils/JDC/bst.t7" | |
ASR_config: "Utils/ASR/config.yml" | |
ASR_path: "Utils/ASR/epoch_00080.pth" | |
PLBERT_dir: 'Utils/PLBERT/' | |
data_params: | |
train_data: "Data/train_list.txt" | |
val_data: "Data/val_list.txt" | |
root_path: "" | |
OOD_data: "Data/OOD_texts.txt" | |
min_length: 50 # sample until texts with this size are obtained for OOD texts | |
preprocess_params: | |
sr: 24000 | |
spect_params: | |
n_fft: 2048 | |
win_length: 1200 | |
hop_length: 300 | |
model_params: | |
multispeaker: true | |
dim_in: 64 | |
hidden_dim: 512 | |
max_conv_dim: 512 | |
n_layer: 3 | |
n_mels: 80 | |
n_token: 178 # number of phoneme tokens | |
max_dur: 50 # maximum duration of a single phoneme | |
style_dim: 128 # style vector size | |
dropout: 0.2 | |
# config for decoder | |
decoder: | |
type: 'hifigan' # either hifigan or istftnet | |
resblock_kernel_sizes: [3,7,11] | |
upsample_rates : [10,5,3,2] | |
upsample_initial_channel: 512 | |
resblock_dilation_sizes: [[1,3,5], [1,3,5], [1,3,5]] | |
upsample_kernel_sizes: [20,10,6,4] | |
# speech language model config | |
slm: | |
model: 'microsoft/wavlm-base-plus' | |
sr: 16000 # sampling rate of SLM | |
hidden: 768 # hidden size of SLM | |
nlayers: 13 # number of layers of SLM | |
initial_channel: 64 # initial channels of SLM discriminator head | |
# style diffusion model config | |
diffusion: | |
embedding_mask_proba: 0.1 | |
# transformer config | |
transformer: | |
num_layers: 3 | |
num_heads: 8 | |
head_features: 64 | |
multiplier: 2 | |
# diffusion distribution config | |
dist: | |
sigma_data: 0.2 # placeholder for estimate_sigma_data set to false | |
estimate_sigma_data: true # estimate sigma_data from the current batch if set to true | |
mean: -3.0 | |
std: 1.0 | |
loss_params: | |
lambda_mel: 5. # mel reconstruction loss | |
lambda_gen: 1. # generator loss | |
lambda_slm: 1. # slm feature matching loss | |
lambda_mono: 1. # monotonic alignment loss (1st stage, TMA) | |
lambda_s2s: 1. # sequence-to-sequence loss (1st stage, TMA) | |
TMA_epoch: 5 # TMA starting epoch (1st stage) | |
lambda_F0: 1. # F0 reconstruction loss (2nd stage) | |
lambda_norm: 1. # norm reconstruction loss (2nd stage) | |
lambda_dur: 1. # duration loss (2nd stage) | |
lambda_ce: 20. # duration predictor probability output CE loss (2nd stage) | |
lambda_sty: 1. # style reconstruction loss (2nd stage) | |
lambda_diff: 1. # score matching loss (2nd stage) | |
diff_epoch: 10 # style diffusion starting epoch (2nd stage) | |
joint_epoch: 15 # joint training starting epoch (2nd stage) | |
optimizer_params: | |
lr: 0.0001 # general learning rate | |
bert_lr: 0.00001 # learning rate for PLBERT | |
ft_lr: 0.00001 # learning rate for acoustic modules | |
slmadv_params: | |
min_len: 400 # minimum length of samples | |
max_len: 500 # maximum length of samples | |
batch_percentage: 0.5 # to prevent out of memory, only use half of the original batch size | |
iter: 20 # update the discriminator every this iterations of generator update | |
thresh: 5 # gradient norm above which the gradient is scaled | |
scale: 0.01 # gradient scaling factor for predictors from SLM discriminators | |
sig: 1.5 # sigma for differentiable duration modeling | |