AI_Demo / AI_Demo.py
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##### `AI_Demo.py`
##### AI Demo, hosted on https://huggingface.co/spaces/DrBenjamin/AI_Demo
##### Please reach out to ben@benbox.org for any questions
#### Loading needed Python libraries
import streamlit as st
import numpy as np
import audio2numpy as a2n
from pydub import AudioSegment
import cv2
from PIL import Image
import torch
from diffusers import StableDiffusionPipeline
from diffusers import StableDiffusionImg2ImgPipeline
from diffusers import StableDiffusionControlNetPipeline, ControlNetModel, UniPCMultistepScheduler
from transformers import pipeline, set_seed
from transformers import VisionEncoderDecoderModel, ViTImageProcessor, AutoTokenizer
import os
os.environ['COMMANDLINE_ARGS'] = '--skip-torch-cuda-test --precision full --no-half'
os.environ['KMP_DUPLICATE_LIB_OK'] = 'True'
#### Functions
### Function predict_step = Image to Text recognition
def predict_step(image):
if image.mode != "RGB":
image = image.convert(mode = "RGB")
pixel_values = feature_extractor(images = image, return_tensors = "pt").pixel_values
pixel_values = pixel_values.to(device)
output_ids = model.generate(pixel_values, **gen_kwargs)
preds = tokenizer.batch_decode(output_ids, skip_special_tokens = True)
preds = [pred.strip() for pred in preds]
return str(preds[0]).capitalize() + '.'
#### Models
st.header('🤗 Hugging Face Diffusers')
st.write('State-of-the-art diffusion models for image, text and audio generation in PyTorch.')
devices = ["mps", "cpu", "cuda"]
device = st.selectbox(label = 'Select device', options = devices, index = 1, disabled = True)
st.write(':orange[MPS for Mac (Metal Performance Shaders), CPU for all systems and CUDA for systems with NVIDIA GPU.]')
models = ["runwayml/stable-diffusion-v1-5", "stabilityai/stable-diffusion-2-1", "hakurei/waifu-diffusion", "stabilityai/stable-diffusion-2-base",
"nlpconnect/vit-gpt2-image-captioning", "openai-gpt", "gpt2-large", "openai/whisper-large-v2"]
model_id_or_path = st.selectbox(label = 'Select model', options = models, index = 5, disabled = True)
if model_id_or_path == "runwayml/stable-diffusion-v1-5":
st.write(':orange[Stable Diffusion v1-5 is the state of the art text-to-image model.]')
elif model_id_or_path == "stabilityai/stable-diffusion-2-1":
st.write(':orange[New stable diffusion text-to-image model at 768x768 resolution.]')
elif model_id_or_path == "stabilityai/stable-diffusion-2-base":
st.write(':orange[New stable diffusion text-to-image model at 512x512 resolution.]')
elif model_id_or_path == "hakurei/waifu-diffusion":
st.write(
':orange[waifu-diffusion is a latent text-to-image diffusion model that has been conditioned on high-quality anime images through fine-tuning.]')
elif model_id_or_path == "nlpconnect/vit-gpt2-image-captioning":
st.write(':orange[vit-gpt2 is an image captioning model.]')
elif model_id_or_path == "openai-gpt":
st.write(
':orange[openai-gpt is a transformer-based language model created and released by OpenAI. The model is a causal (unidirectional) transformer pre-trained using language modeling on a large corpus with long range dependencies.]')
elif model_id_or_path == "gpt2-large":
st.write(
':orange[GPT-2 Large is the 774M parameter version of GPT-2, a transformer-based language model created and released by OpenAI. The model is a pretrained model on English language using a causal language modeling (CLM) objective.]')
elif model_id_or_path == "openai/whisper-large-v2":
st.write(':orange[Whisper is a pre-trained model for automatic speech recognition (ASR) and speech translation.]')
control_net_models = ["None", "lllyasviel/sd-controlnet-canny", "lllyasviel/sd-controlnet-scribble"]
if model_id_or_path == "runwayml/stable-diffusion-v1-5":
disable = False
else:
disable = True
control_net_model = st.selectbox(label = 'Select control net model', options = control_net_models, disabled = disable)
if control_net_model == "lllyasviel/sd-controlnet-canny":
st.write(
':orange[ControlNet is a neural network structure to control diffusion models by adding extra conditions. This checkpoint corresponds to the ControlNet conditioned on Canny edges.]')
elif control_net_model == "lllyasviel/sd-controlnet-scribble":
st.write(
':orange[ControlNet is a neural network structure to control diffusion models by adding extra conditions. This checkpoint corresponds to the ControlNet conditioned on Scribble images.]')
if model_id_or_path != "runwayml/stable-diffusion-v1-5":
control_net_model = "None"
#### Stable diffusion image 2 image with Control Net
if model_id_or_path == "runwayml/stable-diffusion-v1-5" and control_net_model != "None":
with st.form('img2img (Control Net)'):
st.subheader('Image 2 Image (Control Net)')
st.write('Create an image from text input with an image as template.')
image = ''
uploaded_file = st.file_uploader(label = "Upload a picture", type = 'png')
prompt = st.text_input(label = 'Prompt',
value = 'A picture in comic style, bright colours, a house with red bricks, a dark sky with a full yellow moon, best quality, extremely detailed.')
submitted = st.form_submit_button('Submit')
if submitted:
# Check for image data
if uploaded_file is not None:
image = cv2.imdecode(np.frombuffer(uploaded_file.getvalue(), np.uint8), cv2.COLOR_GRAY2BGR)
image = cv2.cvtColor(image, cv2.COLOR_BGR2GRAY)
# Resize image if existend and not 768x640 / 640x768 pixel
h, w = image.shape
if not (h == 768 and w == 640) and not (h == 640 and w == 768):
# Image is bigger in height than width
if h > w:
# Resize cropped image to standard dimensions
image = cv2.resize(image, (640, 768), interpolation = cv2.INTER_AREA)
# Image is smaller in height than width
else:
# Resize cropped image to standard dimensions
image = cv2.resize(image, (768, 640), interpolation = cv2.INTER_AREA)
# Get canny image
image = cv2.Canny(image, 100, 200)
image = image[:, :, None]
image = np.concatenate([image, image, image], axis = 2)
canny_image = Image.fromarray(image)
st.subheader('Preview annotator result')
st.image(canny_image)
# Load control net and stable diffusion v1-5
controlnet = ControlNetModel.from_pretrained(control_net_model, torch_dtype = torch.float32)
pipe = StableDiffusionControlNetPipeline.from_pretrained(model_id_or_path, controlnet = controlnet, torch_dtype = torch.float32)
pipe = pipe.to(device)
# Recommended if your computer has < 64 GB of RAM
pipe.enable_attention_slicing()
# Speed up diffusion process with faster scheduler and memory optimization
pipe.scheduler = UniPCMultistepScheduler.from_config(pipe.scheduler.config)
# Generate image
generator = torch.manual_seed(0)
image = pipe(prompt = prompt, negative_prompt = "monochrome, lowres, bad anatomy, worst quality, low quality", num_inference_steps = 30,
generator = generator, image = canny_image).images[0]
st.subheader('Diffuser result')
st.write('Model :orange[' + model_id_or_path + '] + :red[' + control_net_model + ']')
st.image(image)
## Stable-Diffusion
if model_id_or_path == "runwayml/stable-diffusion-v1-5" and control_net_model == "None":
with st.form('img2img'):
st.subheader('Image 2 Image')
st.write('Create an image from text input with an image as template.')
image = ''
uploaded_file = st.file_uploader(label = "Upload a picture", type = 'png')
prompt = st.text_input(label = 'Prompt',
value = 'A picture in comic style, bright colours, a house with red bricks, a dark sky with a full yellow moon, best quality, extremely detailed.')
submitted = st.form_submit_button('Submit')
if submitted:
# Check for image data
if uploaded_file is not None:
image = cv2.imdecode(np.frombuffer(uploaded_file.getvalue(), np.uint8), cv2.IMREAD_COLOR)
# Resize image if existend and not 768x640 / 640x768 pixel
h, w, _ = image.shape
if not (h == 768 and w == 640) and not (h == 640 and w == 768):
# Image is bigger in height than width
if h > w:
# Resize cropped image to standard dimensions
image = cv2.resize(image, (640, 768), interpolation = cv2.INTER_AREA)
# Image is smaller in height than width
else:
# Resize cropped image to standard dimensions
image = cv2.resize(image, (768, 640), interpolation = cv2.INTER_AREA)
image = Image.fromarray(image)
# Load the pipeline
pipe = StableDiffusionImg2ImgPipeline.from_pretrained(model_id_or_path, torch_dtype = torch.float32)
pipe = pipe.to(device)
# Recommended if your computer has < 64 GB of RAM
pipe.enable_attention_slicing()
# Speed up diffusion process with faster scheduler and memory optimization
pipe.scheduler = UniPCMultistepScheduler.from_config(pipe.scheduler.config)
# Create new image
images = pipe(prompt = prompt, negative_prompt = "monochrome, lowres, bad anatomy, worst quality, low quality", num_inference_steps = 30,
image = image, strength = 0.75, guidance_scale = 7.5).images
# Show image
st.subheader('Diffuser result')
st.write('Model :orange[' + model_id_or_path + ']')
st.image(images[0])
#### Stable diffusion txt 2 image
if control_net_model == "None" and model_id_or_path != "nlpconnect/vit-gpt2-image-captioning" and model_id_or_path != "openai-gpt" and model_id_or_path != "gpt2-large" and model_id_or_path != "openai/whisper-large-v2":
with st.form('txt2img'):
st.subheader('Text 2 Image')
st.write('Create an image from text input.')
if model_id_or_path == "runwayml/stable-diffusion-v1-5" or model_id_or_path == "stabilityai/stable-diffusion-2-1":
value = 'A picture in comic style, bright colours, a house with red bricks, a dark sky with a full yellow moon, best quality, extremely detailed.'
if model_id_or_path == "hakurei/waifu-diffusion":
value = 'A picture in Anime style, bright colours, a house with red bricks, a dark sky with a full yellow moon, best quality, extremely detailed.'
if model_id_or_path == "stabilityai/stable-diffusion-2-base":
value = 'A picture in comic style, a castle with grey bricks in the background, a river is going through, a blue sky with a full yellow sun, best quality, extremely detailed.'
prompt = st.text_input(label = 'Prompt', value = value)
submitted = st.form_submit_button('Submit')
if submitted:
# Make sure you're logged in with `huggingface-cli login`
pipe = StableDiffusionPipeline.from_pretrained(model_id_or_path)
pipe = pipe.to(device)
# Recommended if your computer has < 64 GB of RAM
pipe.enable_attention_slicing()
# Speed up diffusion process with faster scheduler and memory optimization
pipe.scheduler = UniPCMultistepScheduler.from_config(pipe.scheduler.config)
# Results
if model_id_or_path == "hakurei/waifu-diffusion":
negative = "several scenes, more than one image, split picture"
else:
negative = "monochrome, lowres, bad anatomy, worst quality, low quality"
image = pipe(prompt = prompt, negative_prompt = negative, num_inference_steps = 30, guidance_scale = 7.5).images[0]
st.subheader('Diffuser result')
st.write('Model :orange[' + model_id_or_path + ']')
st.image(image)
#### Text (OpenAI gpt models)
if model_id_or_path == "openai-gpt" or model_id_or_path == "gpt2-large":
with st.form('GPT'):
st.subheader('Text generation')
st.write('Create text which is generated from text input.')
text_input = st.text_input(label = 'Give a start of a sentence', value = 'This is a test ')
submitted = st.form_submit_button('Submit')
if submitted:
generator = pipeline('text-generation', model = model_id_or_path)
set_seed(42)
generated = generator(text_input, max_length = 150, num_return_sequences = 1)
st.subheader('Diffuser result')
st.write('Model :orange[' + model_id_or_path + ']')
st.markdown('Text: ":green[' + str(generated[0]['generated_text']) + ']"')
#### Image to text
if model_id_or_path == "nlpconnect/vit-gpt2-image-captioning":
with st.form('Image2Text'):
st.subheader('Image 2 Text')
st.write('Create a description of an image.')
image = ''
uploaded_file = st.file_uploader(label = "Upload a picture", type = 'png')
submitted = st.form_submit_button('Submit')
if submitted:
# Check for image data
if uploaded_file is not None:
image = cv2.imdecode(np.frombuffer(uploaded_file.getvalue(), np.uint8), cv2.IMREAD_COLOR)
image = Image.fromarray(image)
model = VisionEncoderDecoderModel.from_pretrained(model_id_or_path)
feature_extractor = ViTImageProcessor.from_pretrained(model_id_or_path)
tokenizer = AutoTokenizer.from_pretrained(model_id_or_path)
device = torch.device("cuda" if torch.cuda.is_available() else "cpu")
model.to(device)
max_length = 16
num_beams = 4
gen_kwargs = {"max_length": max_length, "num_beams": num_beams}
output = predict_step(image)
st.subheader('Diffuser result')
st.write('Model :orange[nlpconnect/vit-gpt2-image-captioning]')
st.write('Description: ":green[' + str(output) + ']"')
#### Whisper Model
if model_id_or_path == "openai/whisper-large-v2":
with st.form('Image2Text'):
st.subheader('Audio 2 Text')
st.write('Create a transcription of an audio file.')
audio_file = st.file_uploader(label = "Upload an audio file", type = 'mp3')
submitted = st.form_submit_button('Submit')
if submitted:
if audio_file is not None:
audio = audio_file.getvalue()
with open("temp.mp3", "wb") as binary_file:
# Write bytes to file
binary_file.write(audio)
# Calling the split_to_mono method on the stereo audio file
stereo_audio = AudioSegment.from_file("temp.mp3", format = "mp3")
mono_audios = stereo_audio.split_to_mono()
mono_audios[0].export("temp.mp3", format = "mp3")
# Mp3 file to numpy array
audio, sr = a2n.audio_from_file('temp.mp3')
st.audio('temp.mp3')
if os.path.exists("temp.mp3"):
os.remove("temp.mp3")
# Load model and processor
pipe = pipeline("automatic-speech-recognition", model = "openai/whisper-large-v2", chunk_length_s = 30, device = "cpu",
ignore_warning = True)
prediction = pipe(audio, sampling_rate = sr)["text"]
st.subheader('Preview used audio')
st.write('Model :orange[' + model_id_or_path + ']')
st.write('Transcript: ":green[' + str(prediction) + ']"')