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Parent(s): fad99c2
Initial commit
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app.py
CHANGED
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@@ -2,40 +2,35 @@ import gradio as gr
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import torch
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import numpy as np
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import random
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from diffusers import
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from PIL import Image
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dtype = torch.
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device = "cuda" if torch.cuda.is_available() else "cpu"
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MAX_SEED = np.iinfo(np.int32).max
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MODEL_ID = "
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LORA_REPO = "alvdansen/flux-koda"
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LORA_FILE = "flux_dev_koda_araminta_k.safetensors"
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STYLES = {
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"Anime": {
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"
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"file": "flux_dev_koda_araminta_k.safetensors",
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"prompt": "anime style portrait, detailed face, vibrant colors, high quality"
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},
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"Pixel Art": {
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"
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"file": "pixel-art-xl.safetensors",
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"prompt": "pixel art avatar, 16-bit style, retro game character"
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},
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"Ghibli": {
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"
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"file": "flux_ghibsky_illustration.safetensors",
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"prompt": "ghibli studio style illustration, soft colors, dreamy atmosphere"
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},
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}
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print("Chargement du modèle...")
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pipe =
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pipe.to(device)
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current_style = None
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def generate_avatar(
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input_image,
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@@ -45,33 +40,23 @@ def generate_avatar(
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randomize_seed,
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progress=gr.Progress(track_tqdm=True)
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):
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global current_style
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if input_image is None:
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raise gr.Error("Veuillez uploader une photo !")
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# Charger le bon LoRA selon le style
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if style_name != current_style:
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style = STYLES[style_name]
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pipe.unload_lora_weights()
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pipe.load_lora_weights(style["lora"], weight_name=style["file"])
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current_style = style_name
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if randomize_seed:
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seed = random.randint(0, MAX_SEED)
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generator = torch.Generator(device=device).manual_seed(seed)
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prompt = STYLES[style_name]["prompt"]
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# Redimensionner l'image input
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input_image = input_image.resize((512, 512))
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result = pipe(
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prompt=prompt,
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image=input_image,
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strength=strength,
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num_inference_steps=20,
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guidance_scale=
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generator=generator,
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).images[0]
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import torch
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import numpy as np
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import random
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from diffusers import StableDiffusionImg2ImgPipeline
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from PIL import Image
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dtype = torch.float16
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device = "cuda" if torch.cuda.is_available() else "cpu"
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MAX_SEED = np.iinfo(np.int32).max
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MODEL_ID = "runwayml/stable-diffusion-v1-5"
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STYLES = {
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"Anime": {
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"prompt": "anime style portrait, detailed face, vibrant colors, high quality, studio ghibli"
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},
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"Pixel Art": {
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"prompt": "pixel art avatar, 16-bit style, retro game character, sharp pixels"
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},
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"Ghibli": {
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"prompt": "ghibli studio style illustration, soft colors, dreamy atmosphere, watercolor"
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},
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}
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print("Chargement du modèle...")
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pipe = StableDiffusionImg2ImgPipeline.from_pretrained(
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MODEL_ID,
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torch_dtype=dtype,
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safety_checker=None
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)
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pipe.to(device)
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print("Modèle chargé !")
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def generate_avatar(
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input_image,
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randomize_seed,
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progress=gr.Progress(track_tqdm=True)
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):
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if input_image is None:
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raise gr.Error("Veuillez uploader une photo !")
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if randomize_seed:
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seed = random.randint(0, MAX_SEED)
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generator = torch.Generator(device=device).manual_seed(seed)
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prompt = STYLES[style_name]["prompt"]
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input_image = input_image.resize((512, 512))
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result = pipe(
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prompt=prompt,
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image=input_image,
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strength=strength,
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num_inference_steps=20,
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guidance_scale=7.5,
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generator=generator,
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).images[0]
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