nikhilg commited on
Commit
488690a
1 Parent(s): 101a3d0

Update README.md

Browse files
Files changed (1) hide show
  1. README.md +44 -184
README.md CHANGED
@@ -27,37 +27,36 @@ extra_gated_prompt: |-
27
  extra_gated_heading: Please read the LICENSE to access this model
28
  ---
29
 
30
- # Stable Diffusion v1-4 Model Card
31
 
32
- Stable Diffusion is a latent text-to-image diffusion model capable of generating photo-realistic images given any text input.
33
- For more information about how Stable Diffusion functions, please have a look at [🤗's Stable Diffusion with 🧨Diffusers blog](https://huggingface.co/blog/stable_diffusion).
34
 
35
- The **Stable-Diffusion-v1-4** checkpoint was initialized with the weights of the [Stable-Diffusion-v1-2](https:/steps/huggingface.co/CompVis/stable-diffusion-v1-2)
36
- checkpoint and subsequently fine-tuned on 225k steps at resolution 512x512 on "laion-aesthetics v2 5+" and 10% dropping of the text-conditioning to improve [classifier-free guidance sampling](https://arxiv.org/abs/2207.12598).
37
 
38
- This weights here are intended to be used with the 🧨 Diffusers library. If you are looking for the weights to be loaded into the CompVis Stable Diffusion codebase, [come here](https://huggingface.co/CompVis/stable-diffusion-v-1-4-original)
39
 
40
  ## Model Details
41
  - **Developed by:** Robin Rombach, Patrick Esser
42
- - **Model type:** Diffusion-based text-to-image generation model
43
  - **Language(s):** English
44
  - **License:** [The CreativeML OpenRAIL M license](https://huggingface.co/spaces/CompVis/stable-diffusion-license) is an [Open RAIL M license](https://www.licenses.ai/blog/2022/8/18/naming-convention-of-responsible-ai-licenses), adapted from the work that [BigScience](https://bigscience.huggingface.co/) and [the RAIL Initiative](https://www.licenses.ai/) are jointly carrying in the area of responsible AI licensing. See also [the article about the BLOOM Open RAIL license](https://bigscience.huggingface.co/blog/the-bigscience-rail-license) on which our license is based.
45
- - **Model Description:** This is a model that can be used to generate and modify images based on text prompts. It is a [Latent Diffusion Model](https://arxiv.org/abs/2112.10752) that uses a fixed, pretrained text encoder ([CLIP ViT-L/14](https://arxiv.org/abs/2103.00020)) as suggested in the [Imagen paper](https://arxiv.org/abs/2205.11487).
46
- - **Resources for more information:** [GitHub Repository](https://github.com/CompVis/stable-diffusion), [Paper](https://arxiv.org/abs/2112.10752).
47
  - **Cite as:**
48
 
49
- @InProceedings{Rombach_2022_CVPR,
50
- author = {Rombach, Robin and Blattmann, Andreas and Lorenz, Dominik and Esser, Patrick and Ommer, Bj\"orn},
51
- title = {High-Resolution Image Synthesis With Latent Diffusion Models},
52
- booktitle = {Proceedings of the IEEE/CVF Conference on Computer Vision and Pattern Recognition (CVPR)},
53
- month = {June},
54
- year = {2022},
55
- pages = {10684-10695}
56
  }
57
 
 
58
  ## Examples
59
 
60
- We recommend using [🤗's Diffusers library](https://github.com/huggingface/diffusers) to run Stable Diffusion.
61
 
62
  ### PyTorch
63
 
@@ -69,124 +68,34 @@ Running the pipeline with the default PNDM scheduler:
69
 
70
  ```python
71
  import torch
72
- from diffusers import StableDiffusionPipeline
 
73
 
74
- model_id = "CompVis/stable-diffusion-v1-4"
75
  device = "cuda"
76
 
77
-
78
- pipe = StableDiffusionPipeline.from_pretrained(model_id, torch_dtype=torch.float16)
79
  pipe = pipe.to(device)
80
 
81
- prompt = "a photo of an astronaut riding a horse on mars"
82
- image = pipe(prompt).images[0]
83
-
84
- image.save("astronaut_rides_horse.png")
85
- ```
86
-
87
- **Note**:
88
- If you are limited by GPU memory and have less than 4GB of GPU RAM available, please make sure to load the StableDiffusionPipeline in float16 precision instead of the default float32 precision as done above. You can do so by telling diffusers to expect the weights to be in float16 precision:
89
-
90
-
91
- ```py
92
- import torch
93
-
94
- pipe = StableDiffusionPipeline.from_pretrained(model_id, torch_dtype=torch.float16)
95
- pipe = pipe.to(device)
96
- pipe.enable_attention_slicing()
97
-
98
- prompt = "a photo of an astronaut riding a horse on mars"
99
- image = pipe(prompt).images[0]
100
-
101
- image.save("astronaut_rides_horse.png")
102
- ```
103
-
104
- To swap out the noise scheduler, pass it to `from_pretrained`:
105
-
106
- ```python
107
- from diffusers import StableDiffusionPipeline, EulerDiscreteScheduler
108
-
109
- model_id = "CompVis/stable-diffusion-v1-4"
110
-
111
- # Use the Euler scheduler here instead
112
- scheduler = EulerDiscreteScheduler.from_pretrained(model_id, subfolder="scheduler")
113
- pipe = StableDiffusionPipeline.from_pretrained(model_id, scheduler=scheduler, torch_dtype=torch.float16)
114
- pipe = pipe.to("cuda")
115
-
116
- prompt = "a photo of an astronaut riding a horse on mars"
117
- image = pipe(prompt).images[0]
118
-
119
- image.save("astronaut_rides_horse.png")
120
- ```
121
-
122
- ### JAX/Flax
123
-
124
- To use StableDiffusion on TPUs and GPUs for faster inference you can leverage JAX/Flax.
125
-
126
- Running the pipeline with default PNDMScheduler
127
-
128
- ```python
129
- import jax
130
- import numpy as np
131
- from flax.jax_utils import replicate
132
- from flax.training.common_utils import shard
133
-
134
- from diffusers import FlaxStableDiffusionPipeline
135
-
136
- pipeline, params = FlaxStableDiffusionPipeline.from_pretrained(
137
- "CompVis/stable-diffusion-v1-4", revision="flax", dtype=jax.numpy.bfloat16
138
- )
139
-
140
- prompt = "a photo of an astronaut riding a horse on mars"
141
-
142
- prng_seed = jax.random.PRNGKey(0)
143
- num_inference_steps = 50
144
-
145
- num_samples = jax.device_count()
146
- prompt = num_samples * [prompt]
147
- prompt_ids = pipeline.prepare_inputs(prompt)
148
-
149
- # shard inputs and rng
150
- params = replicate(params)
151
- prng_seed = jax.random.split(prng_seed, num_samples)
152
- prompt_ids = shard(prompt_ids)
153
-
154
- images = pipeline(prompt_ids, params, prng_seed, num_inference_steps, jit=True).images
155
- images = pipeline.numpy_to_pil(np.asarray(images.reshape((num_samples,) + images.shape[-3:])))
156
  ```
157
 
158
- **Note**:
159
- If you are limited by TPU memory, please make sure to load the `FlaxStableDiffusionPipeline` in `bfloat16` precision instead of the default `float32` precision as done above. You can do so by telling diffusers to load the weights from "bf16" branch.
160
-
161
- ```python
162
- import jax
163
- import numpy as np
164
- from flax.jax_utils import replicate
165
- from flax.training.common_utils import shard
166
-
167
- from diffusers import FlaxStableDiffusionPipeline
168
-
169
- pipeline, params = FlaxStableDiffusionPipeline.from_pretrained(
170
- "CompVis/stable-diffusion-v1-4", revision="bf16", dtype=jax.numpy.bfloat16
171
- )
172
-
173
- prompt = "a photo of an astronaut riding a horse on mars"
174
-
175
- prng_seed = jax.random.PRNGKey(0)
176
- num_inference_steps = 50
177
-
178
- num_samples = jax.device_count()
179
- prompt = num_samples * [prompt]
180
- prompt_ids = pipeline.prepare_inputs(prompt)
181
-
182
- # shard inputs and rng
183
- params = replicate(params)
184
- prng_seed = jax.random.split(prng_seed, num_samples)
185
- prompt_ids = shard(prompt_ids)
186
-
187
- images = pipeline(prompt_ids, params, prng_seed, num_inference_steps, jit=True).images
188
- images = pipeline.numpy_to_pil(np.asarray(images.reshape((num_samples,) + images.shape[-3:])))
189
- ```
190
 
191
  # Uses
192
 
@@ -203,7 +112,7 @@ tasks include
203
  Excluded uses are described below.
204
 
205
  ### Misuse, Malicious Use, and Out-of-Scope Use
206
- _Note: This section is taken from the [DALLE-MINI model card](https://huggingface.co/dalle-mini/dalle-mini), but applies in the same way to Stable Diffusion v1_.
207
 
208
 
209
  The model should not be used to intentionally create or disseminate images that create hostile or alienating environments for people. This includes generating images that people would foreseeably find disturbing, distressing, or offensive; or content that propagates historical or current stereotypes.
@@ -259,65 +168,16 @@ The concepts are passed into the model with the generated image and compared to
259
 
260
 
261
  ## Training
262
-
263
- **Training Data**
264
- The model developers used the following dataset for training the model:
265
-
266
- - LAION-2B (en) and subsets thereof (see next section)
267
-
268
- **Training Procedure**
269
- Stable Diffusion v1-4 is a latent diffusion model which combines an autoencoder with a diffusion model that is trained in the latent space of the autoencoder. During training,
270
-
271
- - Images are encoded through an encoder, which turns images into latent representations. The autoencoder uses a relative downsampling factor of 8 and maps images of shape H x W x 3 to latents of shape H/f x W/f x 4
272
- - Text prompts are encoded through a ViT-L/14 text-encoder.
273
- - The non-pooled output of the text encoder is fed into the UNet backbone of the latent diffusion model via cross-attention.
274
- - The loss is a reconstruction objective between the noise that was added to the latent and the prediction made by the UNet.
275
-
276
- We currently provide four checkpoints, which were trained as follows.
277
- - [`stable-diffusion-v1-1`](https://huggingface.co/CompVis/stable-diffusion-v1-1): 237,000 steps at resolution `256x256` on [laion2B-en](https://huggingface.co/datasets/laion/laion2B-en).
278
- 194,000 steps at resolution `512x512` on [laion-high-resolution](https://huggingface.co/datasets/laion/laion-high-resolution) (170M examples from LAION-5B with resolution `>= 1024x1024`).
279
- - [`stable-diffusion-v1-2`](https://huggingface.co/CompVis/stable-diffusion-v1-2): Resumed from `stable-diffusion-v1-1`.
280
- 515,000 steps at resolution `512x512` on "laion-improved-aesthetics" (a subset of laion2B-en,
281
- filtered to images with an original size `>= 512x512`, estimated aesthetics score `> 5.0`, and an estimated watermark probability `< 0.5`. The watermark estimate is from the LAION-5B metadata, the aesthetics score is estimated using an [improved aesthetics estimator](https://github.com/christophschuhmann/improved-aesthetic-predictor)).
282
- - [`stable-diffusion-v1-3`](https://huggingface.co/CompVis/stable-diffusion-v1-3): Resumed from `stable-diffusion-v1-2`. 195,000 steps at resolution `512x512` on "laion-improved-aesthetics" and 10 % dropping of the text-conditioning to improve [classifier-free guidance sampling](https://arxiv.org/abs/2207.12598).
283
- - [`stable-diffusion-v1-4`](https://huggingface.co/CompVis/stable-diffusion-v1-4) Resumed from `stable-diffusion-v1-2`.225,000 steps at resolution `512x512` on "laion-aesthetics v2 5+" and 10 % dropping of the text-conditioning to improve [classifier-free guidance sampling](https://arxiv.org/abs/2207.12598).
284
-
285
- - **Hardware:** 32 x 8 x A100 GPUs
286
- - **Optimizer:** AdamW
287
- - **Gradient Accumulations**: 2
288
- - **Batch:** 32 x 8 x 2 x 4 = 2048
289
- - **Learning rate:** warmup to 0.0001 for 10,000 steps and then kept constant
290
-
291
- ## Evaluation Results
292
- Evaluations with different classifier-free guidance scales (1.5, 2.0, 3.0, 4.0,
293
- 5.0, 6.0, 7.0, 8.0) and 50 PLMS sampling
294
- steps show the relative improvements of the checkpoints:
295
-
296
- ![pareto](https://huggingface.co/CompVis/stable-diffusion/resolve/main/v1-variants-scores.jpg)
297
-
298
- Evaluated using 50 PLMS steps and 10000 random prompts from the COCO2017 validation set, evaluated at 512x512 resolution. Not optimized for FID scores.
299
- ## Environmental Impact
300
-
301
- **Stable Diffusion v1** **Estimated Emissions**
302
- Based on that information, we estimate the following CO2 emissions using the [Machine Learning Impact calculator](https://mlco2.github.io/impact#compute) presented in [Lacoste et al. (2019)](https://arxiv.org/abs/1910.09700). The hardware, runtime, cloud provider, and compute region were utilized to estimate the carbon impact.
303
-
304
- - **Hardware Type:** A100 PCIe 40GB
305
- - **Hours used:** 150000
306
- - **Cloud Provider:** AWS
307
- - **Compute Region:** US-east
308
- - **Carbon Emitted (Power consumption x Time x Carbon produced based on location of power grid):** 11250 kg CO2 eq.
309
-
310
 
311
  ## Citation
312
 
313
  ```bibtex
314
- @InProceedings{Rombach_2022_CVPR,
315
- author = {Rombach, Robin and Blattmann, Andreas and Lorenz, Dominik and Esser, Patrick and Ommer, Bj\"orn},
316
- title = {High-Resolution Image Synthesis With Latent Diffusion Models},
317
- booktitle = {Proceedings of the IEEE/CVF Conference on Computer Vision and Pattern Recognition (CVPR)},
318
- month = {June},
319
- year = {2022},
320
- pages = {10684-10695}
321
  }
322
  ```
323
 
 
27
  extra_gated_heading: Please read the LICENSE to access this model
28
  ---
29
 
30
+ # GLIGEN: Open-Set Grounded Text-to-Image Generation
31
 
32
+ The GLIGEN model was created by researchers and engineers from [University of Wisconsin-Madison, Columbia University, and Microsoft](https://github.com/gligen/GLIGEN).
33
+ The [`StableDiffusionGLIGENPipeline`] can generate photorealistic images conditioned on grounding inputs.
34
 
35
+ Along with text and bounding boxes, if input images are given, this pipeline can insert objects described by text at the region defined by bounding boxes.
36
+ Otherwise, it'll generate an image described by the caption/prompt and insert objects described by text at the region defined by bounding boxes. It's trained on COCO2014D and COCO2014CD datasets, and the model uses a frozen CLIP ViT-L/14 text encoder to condition itself on grounding inputs.
37
 
38
+ This weights here are intended to be used with the 🧨 Diffusers library. If you want to use one of the official checkpoints for a task, explore the [gligen](https://huggingface.co/gligen) Hub organizations!
39
 
40
  ## Model Details
41
  - **Developed by:** Robin Rombach, Patrick Esser
42
+ - **Model type:** Diffusion-based Grounded Text-to-image generation model
43
  - **Language(s):** English
44
  - **License:** [The CreativeML OpenRAIL M license](https://huggingface.co/spaces/CompVis/stable-diffusion-license) is an [Open RAIL M license](https://www.licenses.ai/blog/2022/8/18/naming-convention-of-responsible-ai-licenses), adapted from the work that [BigScience](https://bigscience.huggingface.co/) and [the RAIL Initiative](https://www.licenses.ai/) are jointly carrying in the area of responsible AI licensing. See also [the article about the BLOOM Open RAIL license](https://bigscience.huggingface.co/blog/the-bigscience-rail-license) on which our license is based.
45
+ - **Model Description:** This is a model that can be used to generate and modify images based on text prompts and bounding boxes. It is a [Latent Diffusion Model](https://arxiv.org/abs/2112.10752) that uses a fixed, pretrained text encoder ([CLIP ViT-L/14](https://arxiv.org/abs/2103.00020)) as suggested in the [Imagen paper](https://arxiv.org/abs/2205.11487).
46
+ - **Resources for more information:** [GitHub Repository](https://github.com/gligen/GLIGEN), [Paper](https://arxiv.org/pdf/2301.07093.pdf).
47
  - **Cite as:**
48
 
49
+ @article{li2023gligen,
50
+ author = {Li, Yuheng and Liu, Haotian and Wu, Qingyang and Mu, Fangzhou and Yang, Jianwei and Gao, Jianfeng and Li, Chunyuan and Lee, Yong Jae},
51
+ title = {GLIGEN: Open-Set Grounded Text-to-Image Generation},
52
+ publisher = {arXiv:2301.07093},
53
+ year = {2023},
 
 
54
  }
55
 
56
+
57
  ## Examples
58
 
59
+ We recommend using [🤗's Diffusers library](https://github.com/huggingface/diffusers) to run GLIGEN.
60
 
61
  ### PyTorch
62
 
 
68
 
69
  ```python
70
  import torch
71
+ import torchvision
72
+ from diffusers import StableDiffusionGLIGENPipeline
73
 
74
+ model_id = "masterful/gligen-1-4-generation-text-box"
75
  device = "cuda"
76
 
77
+ pipe = StableDiffusionGLIGENPipeline.from_pretrained(model_id, variant="fp16", torch_dtype=torch.float16)
 
78
  pipe = pipe.to(device)
79
 
80
+ prompt = "a water glass and a bread on the kitchen counter"
81
+
82
+ images = pipe(
83
+ prompt,
84
+ num_images_per_prompt=1,
85
+ gligen_phrases = ['a water glass', 'a bread'],
86
+ gligen_boxes = [
87
+ [0.1387, 0.2051, 0.4277, 0.7090],
88
+ [0.4980, 0.4355, 0.8516, 0.7266],
89
+ ],
90
+ gligen_scheduled_sampling_beta=0.3,
91
+ output_type="np",
92
+ num_inference_steps=50
93
+ ).images
94
+
95
+ images = torch.stack([torch.from_numpy(image) for image in images]).permute(0, 3, 1, 2)
96
+ torchvision.utils.save_image(images, "./gligen-1-4-generation-text-box.jpg", nrow=1, normalize=False)
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
97
  ```
98
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
99
 
100
  # Uses
101
 
 
112
  Excluded uses are described below.
113
 
114
  ### Misuse, Malicious Use, and Out-of-Scope Use
115
+ _Note: This section is taken from the [DALLE-MINI model card](https://huggingface.co/dalle-mini/dalle-mini), but applies in the same way to GLIGEN.
116
 
117
 
118
  The model should not be used to intentionally create or disseminate images that create hostile or alienating environments for people. This includes generating images that people would foreseeably find disturbing, distressing, or offensive; or content that propagates historical or current stereotypes.
 
168
 
169
 
170
  ## Training
171
+ Refer [`GLIGEN`](https://github.com/gligen/GLIGEN) for more details.
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
 
172
 
173
  ## Citation
174
 
175
  ```bibtex
176
+ @article{li2023gligen,
177
+ author = {Li, Yuheng and Liu, Haotian and Wu, Qingyang and Mu, Fangzhou and Yang, Jianwei and Gao, Jianfeng and Li, Chunyuan and Lee, Yong Jae},
178
+ title = {GLIGEN: Open-Set Grounded Text-to-Image Generation},
179
+ publisher = {arXiv:2301.07093},
180
+ year = {2023},
 
 
181
  }
182
  ```
183