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Stable Diffusion v2-1 Model Card\nThis model card focuses on the model associated with the Stable Diffusion v2-1 model, codebase available here.\n\nThis stable-diffusion-2-1 model is fine-tuned from stable-diffusion-2 (768-v-ema.ckpt) with an additional 55k steps on the same dataset (with punsafe=0.1), and then fine-tuned for another 155k extra steps with punsafe=0.98.\n\nUse it with the stablediffusion repository: download the v2-1_768-ema-pruned.ckpt here.\nUse it with \uD83E\uDDE8 diffusers\nModel Details\nDeveloped by: Robin Rombach, Patrick Esser\n\nModel type: Diffusion-based text-to-image generation model\n\nLanguage(s): English\n\nLicense: CreativeML Open RAIL++-M License\n\nModel Description: This is a model that can be used to generate and modify images based on text prompts. It is a Latent Diffusion Model that uses a fixed, pretrained text encoder (OpenCLIP-ViT/H).\n\nResources for more information: GitHub Repository.

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Stable Diffusion is a latent text-to-image diffusion model capable of generating photo-realistic images given any text input. For more information about how Stable Diffusion functions, please have a look at \uD83E\uDD17's Stable Diffusion blog.\n\nThe Stable-Diffusion-v1-5 checkpoint was initialized with the weights of the Stable-Diffusion-v1-2 checkpoint and subsequently fine-tuned on 595k steps at resolution 512x512 on \"laion-aesthetics v2 5+\" and 10% dropping of the text-conditioning to improve classifier-free guidance sampling.\n\nYou can use this both with the \uD83E\uDDE8Diffusers library and the RunwayML GitHub repository.

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