# Generative Models by Stability AI ![sample1](assets/000.jpg) ## News **July 4, 2023** - A technical report on SDXL is now available [here](assets/sdxl_report.pdf). **June 22, 2023** - We are releasing two new diffusion models for research purposes: - `SD-XL 0.9-base`: The base model was trained on a variety of aspect ratios on images with resolution 1024^2. The base model uses [OpenCLIP-ViT/G](https://github.com/mlfoundations/open_clip) and [CLIP-ViT/L](https://github.com/openai/CLIP/tree/main) for text encoding whereas the refiner model only uses the OpenCLIP model. - `SD-XL 0.9-refiner`: The refiner has been trained to denoise small noise levels of high quality data and as such is not expected to work as a text-to-image model; instead, it should only be used as an image-to-image model. If you would like to access these models for your research, please apply using one of the following links: [SDXL-0.9-Base model](https://huggingface.co/stabilityai/stable-diffusion-xl-base-0.9), and [SDXL-0.9-Refiner](https://huggingface.co/stabilityai/stable-diffusion-xl-refiner-0.9). This means that you can apply for any of the two links - and if you are granted - you can access both. Please log in to your Hugging Face Account with your organization email to request access. **We plan to do a full release soon (July).** ## The codebase ### General Philosophy Modularity is king. This repo implements a config-driven approach where we build and combine submodules by calling `instantiate_from_config()` on objects defined in yaml configs. See `configs/` for many examples. ### Changelog from the old `ldm` codebase For training, we use [pytorch-lightning](https://www.pytorchlightning.ai/index.html), but it should be easy to use other training wrappers around the base modules. The core diffusion model class (formerly `LatentDiffusion`, now `DiffusionEngine`) has been cleaned up: - No more extensive subclassing! We now handle all types of conditioning inputs (vectors, sequences and spatial conditionings, and all combinations thereof) in a single class: `GeneralConditioner`, see `sgm/modules/encoders/modules.py`. - We separate guiders (such as classifier-free guidance, see `sgm/modules/diffusionmodules/guiders.py`) from the samplers (`sgm/modules/diffusionmodules/sampling.py`), and the samplers are independent of the model. - We adopt the ["denoiser framework"](https://arxiv.org/abs/2206.00364) for both training and inference (most notable change is probably now the option to train continuous time models): * Discrete times models (denoisers) are simply a special case of continuous time models (denoisers); see `sgm/modules/diffusionmodules/denoiser.py`. * The following features are now independent: weighting of the diffusion loss function (`sgm/modules/diffusionmodules/denoiser_weighting.py`), preconditioning of the network (`sgm/modules/diffusionmodules/denoiser_scaling.py`), and sampling of noise levels during training (`sgm/modules/diffusionmodules/sigma_sampling.py`). - Autoencoding models have also been cleaned up. ## Installation: #### 1. Clone the repo ```shell git clone git@github.com:Stability-AI/generative-models.git cd generative-models ``` #### 2. Setting up the virtualenv This is assuming you have navigated to the `generative-models` root after cloning it. **NOTE:** This is tested under `python3.8` and `python3.10`. For other python versions, you might encounter version conflicts. **PyTorch 1.13** ```shell # install required packages from pypi python3 -m venv .pt1 source .pt1/bin/activate pip3 install wheel pip3 install -r requirements_pt13.txt ``` **PyTorch 2.0** ```shell # install required packages from pypi python3 -m venv .pt2 source .pt2/bin/activate pip3 install wheel pip3 install -r requirements_pt2.txt ``` ## Inference: We provide a [streamlit](https://streamlit.io/) demo for text-to-image and image-to-image sampling in `scripts/demo/sampling.py`. The following models are currently supported: - [SD-XL 0.9-base](https://huggingface.co/stabilityai/stable-diffusion-xl-base-0.9) - [SD-XL 0.9-refiner](https://huggingface.co/stabilityai/stable-diffusion-xl-refiner-0.9) - [SD 2.1-512](https://huggingface.co/stabilityai/stable-diffusion-2-1-base/blob/main/v2-1_512-ema-pruned.safetensors) - [SD 2.1-768](https://huggingface.co/stabilityai/stable-diffusion-2-1/blob/main/v2-1_768-ema-pruned.safetensors) **Weights for SDXL**: If you would like to access these models for your research, please apply using one of the following links: [SDXL-0.9-Base model](https://huggingface.co/stabilityai/stable-diffusion-xl-base-0.9), and [SDXL-0.9-Refiner](https://huggingface.co/stabilityai/stable-diffusion-xl-refiner-0.9). This means that you can apply for any of the two links - and if you are granted - you can access both. Please log in to your Hugging Face Account with your organization email to request access. After obtaining the weights, place them into `checkpoints/`. Next, start the demo using ``` streamlit run scripts/demo/sampling.py --server.port ``` ### Invisible Watermark Detection Images generated with our code use the [invisible-watermark](https://github.com/ShieldMnt/invisible-watermark/) library to embed an invisible watermark into the model output. We also provide a script to easily detect that watermark. Please note that this watermark is not the same as in previous Stable Diffusion 1.x/2.x versions. To run the script you need to either have a working installation as above or try an _experimental_ import using only a minimal amount of packages: ```bash python -m venv .detect source .detect/bin/activate pip install "numpy>=1.17" "PyWavelets>=1.1.1" "opencv-python>=4.1.0.25" pip install --no-deps invisible-watermark ``` To run the script you need to have a working installation as above. The script is then useable in the following ways (don't forget to activate your virtual environment beforehand, e.g. `source .pt1/bin/activate`): ```bash # test a single file python scripts/demo/detect.py # test multiple files at once python scripts/demo/detect.py ... # test all files in a specific folder python scripts/demo/detect.py /* ``` ## Training: We are providing example training configs in `configs/example_training`. To launch a training, run ``` python main.py --base configs/ configs/ ``` where configs are merged from left to right (later configs overwrite the same values). This can be used to combine model, training and data configs. However, all of them can also be defined in a single config. For example, to run a class-conditional pixel-based diffusion model training on MNIST, run ```bash python main.py --base configs/example_training/toy/mnist_cond.yaml ``` **NOTE 1:** Using the non-toy-dataset configs `configs/example_training/imagenet-f8_cond.yaml`, `configs/example_training/txt2img-clipl.yaml` and `configs/example_training/txt2img-clipl-legacy-ucg-training.yaml` for training will require edits depending on the used dataset (which is expected to stored in tar-file in the [webdataset-format](https://github.com/webdataset/webdataset)). To find the parts which have to be adapted, search for comments containing `USER:` in the respective config. **NOTE 2:** This repository supports both `pytorch1.13` and `pytorch2`for training generative models. However for autoencoder training as e.g. in `configs/example_training/autoencoder/kl-f4/imagenet-attnfree-logvar.yaml`, only `pytorch1.13` is supported. **NOTE 3:** Training latent generative models (as e.g. in `configs/example_training/imagenet-f8_cond.yaml`) requires retrieving the checkpoint from [Hugging Face](https://huggingface.co/stabilityai/sdxl-vae/tree/main) and replacing the `CKPT_PATH` placeholder in [this line](configs/example_training/imagenet-f8_cond.yaml#81). The same is to be done for the provided text-to-image configs. ### Building New Diffusion Models #### Conditioner The `GeneralConditioner` is configured through the `conditioner_config`. Its only attribute is `emb_models`, a list of different embedders (all inherited from `AbstractEmbModel`) that are used to condition the generative model. All embedders should define whether or not they are trainable (`is_trainable`, default `False`), a classifier-free guidance dropout rate is used (`ucg_rate`, default `0`), and an input key (`input_key`), for example, `txt` for text-conditioning or `cls` for class-conditioning. When computing conditionings, the embedder will get `batch[input_key]` as input. We currently support two to four dimensional conditionings and conditionings of different embedders are concatenated appropriately. Note that the order of the embedders in the `conditioner_config` is important. #### Network The neural network is set through the `network_config`. This used to be called `unet_config`, which is not general enough as we plan to experiment with transformer-based diffusion backbones. #### Loss The loss is configured through `loss_config`. For standard diffusion model training, you will have to set `sigma_sampler_config`. #### Sampler config As discussed above, the sampler is independent of the model. In the `sampler_config`, we set the type of numerical solver, number of steps, type of discretization, as well as, for example, guidance wrappers for classifier-free guidance. ### Dataset Handling For large scale training we recommend using the data pipelines from our [data pipelines](https://github.com/Stability-AI/datapipelines) project. The project is contained in the requirement and automatically included when following the steps from the [Installation section](#installation). Small map-style datasets should be defined here in the repository (e.g., MNIST, CIFAR-10, ...), and return a dict of data keys/values, e.g., ```python example = {"jpg": x, # this is a tensor -1...1 chw "txt": "a beautiful image"} ``` where we expect images in -1...1, channel-first format.